Pine Waves 2015

Oscillations & Waves
An intermediate introduction to the physics of waves
David J. Pine
and
T. C. Lubensky
c Copyright 2015

D. J. Pine & T. C. Lubensky
September 27, 2015
Contents
List of illustrations page viii

List of tables xi
Part I Oscillations 1
1 Oscillations 3
1.1 The simple pendulum 3
1.2 Diatomic molecule 7
1.3 The ideal spring 10
1.3.1 The horizontal spring 11
1.3.2 The vertical spring: an inhomogeneous ODE 11
1.3.3 Initial conditions and constraints 14
1.4 Energy of a simple harmonic oscillator 15
1.5 Complex exponential notation 19
1.5.1 Complex exponentials 19
1.5.2 Complex exponentials in harmonic oscillator problems 21
1.6 Series expansions 23
1.7 Summary of important points of Chapter 1 25
Problems 27
2 Damped Oscillations 31
2.1 The damped oscillator 31
2.1.1 Underdamped response: (γ/2)2 < ω20 34
2.1.2 Overdamped response: (γ/2)2 > ω20 37
2.1.3 Critically damped oscillations: (γ/2)2 = ω20 39
2.1.4 Summary of damped oscillations 40
2.2 Energy loss from damping 41
2.3 Shock absorbers 43
Problems 48
3 Resonance 52
3.1 Forced oscillators and resonance 52
3.1.1 Resonance 53
3.1.2 Response functions 59
3.1.3 Dissipation 61
iii
iv Contents
t
3.1.4 Energy stored in a driven oscillator 62
3.1.5 Transients 63
3.1.6 Indirect forcing 65
3.2 Resonant filters and detectors 68
3.2.1 Radio receiver 70
Problems 74
4 Normal Modes 78
4.1 Systems with two degrees of freedom 78
4.1.1 Two coupled pendulums 78
4.1.2 Normal modes of two coupled pendulums 82
4.1.3 Weak coupling 86
4.1.4 Energy in normal modes 88
4.2 Matrix formulation of normal modes 89
4.2.1 Two equal-mass particles 90
4.2.2 The meaning of normal coordinates 94
4.2.3 Specific solutions 97
4.2.4 Three masses on a string 98
4.2.5 Systems of unequal masses 103
4.2.6 Geometry and symmetry 107
4.3 Normal modes of carbon dioxide 109
4.3.1 Longitudinal modes 109
4.3.2 Transverse modes 113
4.4 Damping and normal modes 116
4.5 Forced oscillations 119
4.5.1 Steady-state response 119
4.6 Summary of important points of Chapter 4 122
Problems 123
5 Waves of oscillating particles 129

5.1 Identical point masses on a massless string 129
5.1.1 Small N solutions 131
5.1.2 Extrapolation to large N solutions 133
5.2 Identical masses coupled by springs 140
5.2.1 Masses coupled by springs with fixed ends 140
5.2.2 Periodic boundary conditions 141
5.3 Alternating masses 144
5.3.1 Eigenfrequencies of alternating masses on a string 146
5.3.2 Eigenvectors for periodic boundary conditions 147
5.3.3 Eigenvectors for fixed boundary conditions 151
5.4 Nonperiodic systems 154
5.4.1 Diagonalizing the dynamical matrix 154
5.4.2 Random sequences of masses: Localization 154
5.4.3 A Fibonacci sequence of two masses: Quasiperiodicity 161
v Contents
t
Problems 164
6 Fourier Analysis 166

6.1 Fourier series 167
6.1.1 Sine and cosine Fourier series 167
6.1.2 Fourier series in time 172
6.1.3 Basis functions and orthogonality 172
6.1.4 Power in periodic waves 173
6.1.5 Complex exponential Fourier series 175
6.2 Fourier transforms 176
6.2.1 Fourier transform pairs 176
6.2.2 Some examples 179
6.2.3 The uncertainty principle 181
6.2.4 Some theorems 183
6.3 Discrete Fourier Transforms 184
Problems 190
7 Strings 196
7.1 Waves on strings 196
7.2 Normal modes of a string 198
7.2.1 String with fixed ends 200
7.2.2 String with periodic boundary conditions 203
7.2.3 Traveling waves 203
7.2.4 Standing waves 205
7.3 General solutions of the wave equation 206
7.4 Pulses on strings 207
7.5 Energy of a vibrating string 208
7.5.1 Energy in a pulse 209
7.5.2 Energy density of a traveling pulse 210
7.5.3 Transport of energy 212
7.5.4 Energy density of a traveling wave 214
7.5.5 Energy of standing waves 214
7.6 Boundaries 215
Problems 216
8 Sound and Electromagnetic Waves 217

8.1 Sound 217
8.2 Reflection and refraction 217
8.3 Electromagnetic waves 217
Problems 218
9 Interference, Diffraction, and Fourier Optics 219

9.1 Interference 219
9.2 Diffraction 219
vi Contents
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9.3 Fourier optics 219
9.4 Scattering 219
Problems 220
vii Contents
t
Part II Closing features 221
Appendix A Linear algebra 223

A.1 Vectors 223
A.2 Matrices 224
A.2.1 Matrix operations 224
A.3 Properties of matrices 225
A.3.1 Determinant of a matrix 225
A.3.2 Inverse of a matrix 226
A.3.3 Matrix symmetry 226
A.3.4 Hermitian matrices 227
A.3.5 Positive definite matrices 227
A.3.6 Orthogonal matrices 227
A.3.7 Matrix identities 227
A.4 Transformations 228
A.4.1 Rotations 228
A.4.2 Choleski decomposition 228
A.4.3 Similarity transformations 230
A.5 Eigenvalue equations 230
Appendix B Lagrangian formulation of classical mechanics 232
Appendix C Computer programs 242
Notes 243
Illustrations
1.1 Simple pendulum 4

1.2 Sinusoidal trajectory 5
1.3 Approximations for sin θ. 6
1.4 Potential energy for a diatomic molecule 8
1.5 Ideal spring on a frictionless surface 10
1.6 Vertical spring 12
1.7 Mass with spring on inclined plane 13
1.8 Colliding masses with springs 15
1.9 Energies for mass on a spring as a function of position 17
1.10 Stretched springs 18
1.11 The complex plane 20
1.12 Oscillator on Inclined Plane 29
1.13 Four Springs 30
1.14 Drooping sprints 30
2.1 Damped pendulum 32
2.2 Transient response of an underdamped oscillator 36
2.3 Transient response of an overdamped oscillator 38
2.4 Transient response for damped oscillators 40
2.5 Automobile suspension system with shock absorber 44
2.6 Response of automobile shock absorbers 46
2.7 Torsion Oscillator 50
3.1 Damped forced mass suspended from a spring 53
3.2 Amplitude and phase of sinusoidally driven oscillator 56
3.3 Trajectories of oscillators driven at and off resonance 58
3.4 Response functions of an oscillator 59
3.5 Resonant power absorption of an oscillator 61
3.6 Transient response of a driven oscillator 66
3.7 Pendulum with a moving pivot point 67
3.8 Amplitude and phase of shaking pendulum 69
3.9 Amplitude vs. frequency for seismometer pendulum 70
3.10 Filtered and unfiltered radio signal 70
3.11 Resonant circuit with a resistor, inductor, and capacitor in series. 72
3.12 Resonant RLC circuit. 76
4.1 Two identical pendulums coupled by a spring 79
4.2 Modes of two pendulums coupled by a spring. 83
4.3 Displacements of two strongly coupled pendulums 86
viii
ix Illustrations
t
4.4 Displacements of two weakly coupled pendulums 87
4.5 Horizontal masses on a frictionless surface coupled by two springs 90
4.6 Normal coordinates for 2-mass 2-spring problem. 95
4.7 Three masses on a string under tension 98
4.8 Normal modes of 3 masses on a string 101
4.9 Trajectories of weakly-coupled pendulums of unequal masses 107
4.10 Normal mode coordinate transformation 108
4.11 Longitudinal and transverse vibrations of CO2 109
4.12 Normal modes of CO2 112
4.13 Transverse modes of CO2 : linear and circular polarization 115
4.14 Resonance curves for two weakly coupled pendulums with forcing 122
4.15 Wilberforce pendulum 126
5.1 N masses on a string 130
5.2 Normal modes for N masses on a string for N = 1, 2, 3. 132
5.3 Normal frequencies vs. mode number for identical masses on a string 135
5.4 Normal modes for five particles on a string 137
5.5 N spring-coupled masses 141
5.6 N spring-coupled masses 141
5.7 Modes of N spring-coupled masses with periodic boundary conditions 143
5.8 Sequence of alternating masses on a string under tension 145
5.9 Spectrum of normal frequencies for a repeating sequence of unequal masses 147
5.10 Eigenvectors for alternating masses with periodic boundary conditions 149
5.11 Coordinates for alternating masses on a string 152
5.12 Eigenvectors for alternating masses with fixed boundary conditions 153
5.13 Normal modes for a random sequence of particles 156
5.14 Normal modes for a different random sequence of particles 157
5.15 High frequency normal modes for a random sequence of particles 157
5.16 Localization length vs. frequency 158
5.17 Normal modes for a random sequence of 32 spheres 159
5.18 Density of states for a random sequence of masses 160
5.19 Frequencies of normal modes for a Fibonacci sequence 162
5.20 Normal modes for a Fibonacci sequence of two masses 162
5.21 Normal modes for an 89-mass Fibonacci sequence of two masses 163
6.1 Periodic sawtooth and triangle waveforms 167
6.2 Cosine and sine waves for the first three terms in a Fourier series 168
6.3 Fourier series of sawtooth function 170
6.4 Fourier components of sawtooth function 171
6.5 Sinusoidal wave packet with Gaussian envelope 177
6.6 Fourier transform of sinusoidal wave packet with Gaussian envelope 179
6.7 (a) Square pulse and (b) its Fourier transform. 181
6.8 Discrete sampling for discrete Fourier transform 185
6.9 Discrete Fourier transform of an exponential wave packet 188
6.10 Filtering a noisy Gaussian pulse 190
7.1 Discrete masses to a string with continuous mass distribution 197
x Illustrations
t
7.2 Normal modes of a string 201
7.3 Traveling wave 204
7.4 Standing wave from counter-propagating traveling waves 206
7.5 Traveling pulses 208
7.6 Traveling Gaussian pulse 209
7.7 Geometry for calculating potential energy of a string segment 210
B.1 Double pendulum. 237
Tables
6.1 Ordering of spatial frequencies for the discrete Fourier transform. 187
xi
PART I
OSCILLATIONS
1 Oscillations
Everything oscillates: stars revolving in galaxies, planets orbiting stars, seas, clocks, hearts—
the list is endless. The space all around you is filled with radio waves that oscillate. There
are theories that propose that the universe itself oscillates. Even when you think you’ve
found something that doesn’t oscillate—say the chair you are sitting on—you find that it
is made of molecules and atoms that oscillate, and the atoms themselves are made up other
things that oscillate: electrons, protons, and neutrons. When we look even deeper into what
makes up these tiny constituents, we find more things that oscillate: photons, quarks, and
maybe even strings. So if our goal is to understand the world around us, the world within
us, and the worlds beyond us, we need to understand oscillations.
In this chapter, we develop the kinematics and dynamics of sinusoidal oscillations. We
do so by first considering two physical examples: an oscillating simple pendulum and a
vibrating diatomic molecule. We then turn to a common example, a mass attached to an
ideal spring. We also begin what will be an ongoing task of developing mathematical tools
for describing oscillations. In this chapter, we introduce complex exponentials, a wonderful
mathematical tool that greatly simplifies the algebra of sinusoidal functions.
1.1 The simple pendulum
We begin our discussion of oscillations with the simple pendulum, a mass hung from a fixed
point so that it can freely swing back and forth. When pulled from its equilibrium position
and released, gravity causes the pendulum to fall and the pendulum begins oscillating. In
reality, the oscillations die out in time because of friction, but in this chapter we consider
only the idealized case where there is no friction. In Chapter 2, we introduce friction.
Our first task is to describe the motion of the pendulum quantitatively. To this end, we
reduce the pendulum to its simplest form, a mass m attached to one end of a rigid massless
rod of length l that pivots about its other end, as shown in Fig. 1.1. When the pendulum is
displaced from θ = 0, the rod is constrained to move along an arc described by s = lθ. The
component of gravity directed along the arc s provides a restoring force Fr back towards
θ = 0. It is given by
Fr = −mg sin θ . (1.1)
To obtain the equation of motion for the pendulum, we apply Newton’s second law along
3
4 Oscillations
t
t
Fig. 1.1 Simple pendulum, captured at an instant when it is moving towards θ = 0.
the arc of the pendulum’s path:

d2 s
m = −mg sin θ . (1.2)
dt2
Using s = lθ to eliminate s, Eq. (1.2) yields a differential equation for the angular displace-
ment θ of the pendulum
d2 θ g
+ sin θ = 0 . (1.3)
dt2 l
Equation (1.3) is not easily solved for arbitrary values of θ. However, in the limit of small
oscillations, that is, when θ 1, we can expand sin θ in a Taylor series,
1 3 1 5
sin θ = θ − 3! θ + 5! θ − ... , (1.4)
and obtain an analytic solution to Eq. (1.3) by keeping only the first term. Using the ap-
proximation sin θ ' θ, Eq. (1.3) becomes
d2 θ
+ ω20 θ = 0 , (1.5)
dt2
where we have defined the frequency ω0 ≡ g/l. You can readily verify that the units for
p
ω0 are [time]−1 , as g has units of acceleration [length]/[time2 ] and l has units of [length].
We pause briefly to make some general mathematical observations about Eq. (1.5). We
have written Eq. (1.5) so that all the terms involving the dependent variable θ appear on
the left hand side; any terms not involving θ should appear on the right hand side. This is
a standard form for writing differential equations. Equation (1.5) is a homogeneous linear
second order ordinary differential equation (ODE). It is homogeneous because the right
hand side is zero; there is no term in the equation that does not involve the dependent
variable θ. It is second order because the highest derivative appearing in the equation is a
second derivative. It is ordinary because it involves only one independent variable, t in this
case. It is linear because the dependent variable θ appears only as θ or as derivatives of θ;
there are no terms like θµ where µ , 1 or terms like θµ θ(n) where µ , 0 and θ(n) denotes
5 The simple pendulum
t
θ (t)
Ω0 = ω 0 b
A θ0 = a
T
t/T
1 2
tδ
t
Fig. 1.2 Sinusoidal trajectory of simple harmonic motion described by Eq. (1.6). The initial
displacement is a and the initial (angular) velocity is ω0 b. The period is T = 2π/ω0 .
the nth derivative (n ≥ 1) of θ. We will make use of these mathematical properties, noting
when we do so, in this and the following chapters.
You probably know from your previous studies that the solutions to Eq. (1.5) are si-
nusoidal. In fact both sine and cosine functions are solutions. That there are two linear
independent solutions, sine and cosine, follows from the fact that Eq. (1.5) is a second
order linear ODE; third order linear ODEs have three linearly independent solutions, ...
Because Eq. (1.5) is linear, the sum of of any set of solutions is also a solution. This is the
principle of linear superposition. In particular, any sum of the sine and cosine solutions is
also a solution. Thus, the most general solution is
θ(t) = a cos ω0 t + b sin ω0 t . (1.6)
The solutions satisfy Eq. (1.5) only when the coefficient of t in the sine and cosine functions
is ω0 . You can check this by substituting Eq. (1.6) into Eq. (1.5). The cosine and sine
functions are periodic with period 2π, which means that the solution θ(t) is periodic with
period T = 2π/ω0 . Mathematically, we can express the periodicity of θ(t) by writing θ(t) =
θ(t + T ). The angular frequency ω0 has units of radians per second and should not be
confused with the cyclic frequency ν0 = 1/T measured in cycles per second. The two
frequencies are related to each other by ω0 = 2πν0 . As 2π has no units, both ω0 and ν0
have the same units of inverse time; radians and cycles are unitless.
The constants a and b are determined by the initial conditions. Setting t = 0, we obtain
the following equations for the initial position and velocity from Eq. (1.6)
θ(0) ≡ θ0 = a cos(ω0 · 0) + b sin(ω0 · 0) = a (1.7)

θ̇(0) ≡ Ω0 = −a ω0 sin(ω0 · 0) + b ω0 cos(ω0 · 0) = ω0 b , (1.8)
6 Oscillations
t
y
y=θ y = sin θ
1.0
1 3
y = θ − 3! θ
0.5
0.0 θ
0.0 0.5 1.0 1.5 2.0
t
Fig. 1.3 Approximations for sin θ.
where we have introduced the notation θ̇ ≡ dθ/dt = Ω(t). Thus

Ω0
θ(t) = θ0 cos ω0 t + sin ω0 t. (1.9)
ω0
This equation can be expressed in a different but equivalent form that is often useful:
θ(t) = A cos(ω0 t − δ) (1.10)

= A(cos ω0 t cos δ + sin ω0 t sin δ) . (1.11)
Equating the coefficients of cos ω0 t and sin ω0 t in Eq. (1.9) and Eq. (1.10) leads to
Ω0
q
A = θ02 + (Ω0 /ω0 )2 ; δ = tan−1 . (1.12)
ω0 θ0
Figure 1.2 shows a typical function θ(t) in which both θ0 and Ω0 , which is the slope of θ(t)
at t = 0, are positive. It is clear from Eq. (1.10) that θ(t) reaches its maximum value of A at
time tδ = δ/ω0 = δT/2π.
The motion of an oscillator whose equation of motion has the same form as Eq. (1.5)
is often called simple harmonic motion because its trajectory is described by a simple
sinusoidal (harmonic) function. From the point of view of mathematics, it is certainly the
simplest kind of oscillator to treat. It is important to keep in mind, however, that to obtain
this simple result, we had to replace the true restoring force given by Eq. (1.1) with the
linearized approximate force
Fr` = −mg θ . (1.13)
Therefore, our treatment is an approximation that is valid only when sin θ ' θ, that is, when
θ 1. We could obtain more accurate results by including more terms from Eq. (1.4) in
our approximation of the force law:
Fr ' −mg θ − 3!1 θ3 + 5!1 θ5 − ...

(1.14)
= −mg θ 1 − 3!1 θ2 + 5!1 θ4 − ... .

(1.15)
7 Diatomic molecule
t
When a force law contains terms that are nonlinear in θ, the force is said to be anharmonic.
In this case, the lowest order anharmonic term is cubic. The anharmonic terms can be ig-
nored to the extent that they are small compared to the linear, or harmonic, term. Examining
Fig. 1.3, it appears that the linear approximation is good out to about θ ≈ 0.5 rad. Including
the first anharmonic term gives an adequate approximation out to about θ ≈ 1 rad. If we
limit the amplitudes of oscillation to be less than 0.5 rad or about 30◦ , the relative error in
the force incurred by keeping only the linear term is approximately
1 2 1 4
h i
∆F Fr − Fr` −mg θ 1 − 3! θ + 5! θ − ... − 1
≡ = ' 16 θ2 ' 0.04 , (1.16)
Fr Fr` −mgθ
or about 4%. Equation (1.16) also implies that reducing the amplitude of the pendulum by
another factor of two reduces the error in the force by a factor of four. So the harmonic
approximation can be made as accurate as we please provided we limit our description to
sufficiently small amplitudes of the pendulum.
In the analysis of oscillating systems, one rarely encounters a purely linear force law
for displacements from the equilibrium position. However, expanding the force law in a
Taylor series and keeping only the first term, as we did above for the case of the pendulum,
is almost always the first line of attack in trying to describe an oscillating system. Doing
so allows us to solve the equation of motion analytically. Including higher order terms in
the Taylor series expansion of sin θ requires analytical approximation methods or numer-
ical techniques. The range of displacements over which such an approximation is valid
depends on the force law that describes the physical system and the accuracy required of
the solution. In any real system, you should always ask yourself how much error is incurred
by keeping only the linear term by performing an analysis similar to that done above. In
the next section, we consider molecular vibrations, another system for which the harmonic
approximation is often made.
1.2 Diatomic molecule
The atoms in a diatomic molecule vibrate back and forth, alternately stretching and com-
pressing the bond between atoms, as the system exchanges kinetic and potential energy.
The interatomic potential energy is complicated but is sometimes approximated by the
Morse potential, a function given by
2
U(r) = U0 1 − e−(r−r0 )/σ ,

(1.17)
where r = R2 − R1 is the distance between atoms, as illustrated in Fig. 1.4. The potential
energy well has a minimum at r = r0 , the equilibrium distance between the atoms, and has
a width of approximately σ. Figure 1.4 shows a plot of Eq. (1.17) for the case σ = 0.1r0 .
For r < r0 , which corresponds to the two atoms being closer than the equilibrium distance,
the potential is strongly repulsive due to the electron clouds around each atom. For r r0 ,
the potential increases from its minimum of zero at r = r0 towards a value of U0 . This is the
8 Oscillations
t
U0
r
R
R1 R2
1
2 U0
r
R
R1 R2
0 r/ r0
t
0.9 1.0 1.1 1.2 1.3
Fig. 1.4 Potential energy for a diatomic molecule based on Eq. (1.17) with σ = 0.1r0 . The solid
trace shows the potential energy U(r) between two covalently bound atoms as a
function of the distance r ≡ R2 − R1 between them. The potential is a minimum at the
equilibrium separation r0 . The dashed line shows the harmonic approximation U s (δr)
where δr ≡ r − r0 . The coordinates of the atomic nuclei are shown on the right
together with a crude indication of the electron clouds. The two atoms, located at
R = R1 and R = R2 , are shown in their compressed and stretched configurations in the
upper and lower diagrams, respectively
binding energy between the two atoms; the atoms become unbound if the kinetic energy
associated with the vibrational motion of the atoms exceeds this value.
When the separation between the two atoms in the molecule deviates from r = r0 , the
atoms oscillate back and forth. Analyzing the oscillations about the equilibrium position
at r = r0 leads to a nonlinear differential equation that cannot be solved analytically if we
use the full potential given by Eq. (1.17). To simplify the analysis, the potential can be
expanded in a Taylor series about r = r0 where the potential is a minimum (zero in this
case):
U 00 (r0 ) U (n) (r0 )
U(r) = U(r0 ) + U 0 (r0 )(r − r0 ) + (r − r0 )2 + ... + (r − r0 )n + ... (1.18)
2! n!
Applying this formula to Eq. (1.17) while noting that U(r0 ) = 0 and U 0 (r0 ) = 0 gives
" 2
δr δr3 7 δr4
#
U(r) = U0 2 − 3 + + ... , (1.19)
σ σ 12 σ4
where δr ≡ r − r0 is the displacement from the equilibrium distance between the atoms.
Keeping only the first term in the series expansion results in a potential energy that is
quadratic in the displacement δr. In this case, the potential energy can be written in a form
equivalent to that of an ideal spring,
1
k δr2 .
U s (δr) = (1.20)
2
Comparing Eq. (1.20) to Eq. (1.19) we see that the spring constant is given by
k = 2U0 /σ2 . (1.21)
9 Diatomic molecule
t
The dashed line in Fig. 1.4 shows U s (δr). It is apparent from the plot that the approximation
is accurate only if the oscillations are quite small, something like a few per cent of the
distance r0 between the atoms. You can verify that the units of k are those of a spring,
[force]/[length], as U0 is an energy with units of [force][length] (i.e. work) and δr has units
of [length].
The force law corresponding to the potential given by Eq. (1.21) can be obtained from
the gradient (spatial derivative) of the potential:
∂U s (r)
Fs = − = −k δr . (1.22)
∂r
We recognize Eq. (1.22) as the usual linear spring force law where the force is proportional
to minus the displacement from the equilibrium position. This force law has the same form
as Eq. (1.13), the linearized force law for the pendulum. In both cases, the restoring force
(note the minus sign) is proportional to the displacement. For the pendulum the “spring
constant" is mg/l, while for the oxygen molecule it is 2U0 /σ2 . Because the force laws for
the pendulum and the diatomic molecule are the same in the limit of small displacements,
we expect their behavior to be the same, at least for small amplitudes.
To find explicit equations for the motion of the atoms, we need to write down Newton’s
second law. Let’s assume that the two atoms are the same element, say oxygen, so both
atoms have the same mass m. Referring to Fig. 1.4 for the coordinates, an equation of
motion can be written for each particle
d 2 R1
F1 = m = k(R2 − R1 − r0 ) (1.23)
dt2
d 2 R2
F2 = m 2 = −k(R2 − R1 − r0 ) (1.24)
dt
We are looking for an equation that describes the vibration of the two atoms, that is, an
equation that governs how the distance R2 − R1 between the two atoms oscillates in time.
To this end, we subtract Eq. (1.23) from Eq. (1.24), which gives
d2 (R2 − R1 )
m = −2k(R2 − R1 − r0 ) . (1.25)
dt2
Instead of working with R2 − R1 , we subtract off the equilibrium separation r0 between the
two atoms and define δr = R2 − R1 − r0 , which at equilibrium is zero. With this definition
Eq. (1.25) becomes
d2 δr
µ = −k δr , (1.26)
dt2
where we have made use of the fact that r0 is a constant so its derivative is zero. In obtaining
Eq. (1.26), we have reduced the two-body problem (with two masses and two equations)
to a one-body problem with an effective masspof µ. The constant µ ≡ m/2 and is called the
reduced mass of the system. Defining ω0 = k/µ, Eq. (1.26) can be rewritten as
d2 δr
+ ω20 δr = 0 , (1.27)
dt2
which has the same form as Eq. (1.5) and therefore has the same solution with a sinusoidal
10 Oscillations
t
form given by Eq. (1.6). This means that the diatomic atom oscillates with a frequency
ω0 = k/µ, where k is given by Eq. (1.21).
p
Our analysis shows that approaching the problem of oscillations starting from conserva-
tion of energy yields the same results as does approaching the problem starting from forces
and Newton’s second law. Evidently, keeping only terms up to the quadratic term in the
energy formulation is equivalent to keeping the linear term in the force formulation.
An important point to take away from the analysis is that physical systems that can be
described by the same differential equations exhibit the same physical behavior. This is a
very powerful idea because it means that once you understand one system, for example
the pendulum as a harmonic oscillator, you understand the essential physics of many sys-
tems. Here we discover that diatomic molecules also behave like harmonic oscillators if
the displacement from the equilibrium position is sufficiently small. In fact, as we shall
see, nearly all vibrating systems behave like harmonic oscillators provided the amplitude
of the oscillations is sufficiently small.
Exercise 1.2.1 Over what range of distances is the quadratic approximation for Eq.
(1.17) valid to within 10% of the exact version? Use parameters for O2 : bond length
r0 = 0.121 nm, σ = 0.1r0 , and bond energy U0 = 5.2 eV. Hint: Using the Taylor
expansion of U(r) given by Eq. (1.19), find an equation for δU(r)/U(r) ≡ [U(r) −
U s (r)]/(U(r).
1.3 The ideal spring
Nothing in life is perfect, so the saying goes, yet physicists are forever analyzing the ideal
spring. As in the rest of life, real springs do not achieve our ideal. In particular, they do not
obey a simple linear force law, except to some degree of approximation that is valid for suf-
ficiently small displacements. In this way, real springs are just like pendulums and diatomic
molecules. Nevertheless, the ideal spring serves a launching point for most discussions of
oscillations.
Our purpose here in not having started our discussion of oscillations with the ideal spring
is to emphasize that the harmonic approximation is just that, an approximation that is
t
Fig. 1.5 Ideal spring on a frictionless surface for unstretched spring (top) and stretched sprint
(bottom).
11 The ideal spring
t
generally valid only over some finite range of displacements from the equilibrium position.
With that fact firmly in mind, we briefly reformulate the problem of small oscillations in
terms of the ideal spring and recapitulate the main results of the previous sections.
1.3.1 The horizontal spring
We adopt the notation usually associated with an ideal spring with the familiar linear force
law (Hooke’s Law)
F = −kx , (1.28)
where x is the distance the spring is stretched (see Fig. 1.5). The spring is attached to a
mass m that moves horizontally on a frictionless surface, as shown in Fig. 1.5 (recall that
we are treating only frictionless cases in this chapter). Applying Newton’s second law to
the system, we obtain the equation of motion
F = −kx = m ẍ , (1.29)
which can be rewritten in the familiar form
d2 x
+ ω20 x = 0 , (1.30)
dt2
√
where ω0 = k/m. Equation (1.30) has the same form and the same solutions as the
linearized equation of motion Eq. (1.5) for the pendulum. In particular, the general solution
to Eq. (1.30) can be written as
x(t) = a cos ω0 t + b sin ω0 t , (1.31)
or equivalently as
x(t) = d cos(ω0 t − δ); , (1.32)
where a, b, d, and φ can be determined from the initial conditions.
1.3.2 The vertical spring: an inhomogeneous ODE
Suppose a mass hangs vertically from a spring attached to the ceiling as shown in Fig. 1.6.
If the length of the unstretched spring is L0 and y = L − L0 is the distance the spring is
stretched downward from this position, then the spring force −ky on the mass is upward
and negative, while the gravity force mg is downward and positive. The equation of motion
is
d2 y
m = −ky + mg . (1.33)
dt2
Rearranging terms a bit, we obtain
d2 y k
+ y=g. (1.34)
dt2 m
√
Setting ω0 = k/m, the left hand side of this equation is the same equation as Eq. (1.30)
above, but now it is set equal to a constant term rather than zero on the right hand side.
12 Oscillations
t
t
Fig. 1.6 (a) Displacement of a mass hanging from a vertical spring for the unstretched (left)
and stretched positions (right). (b) Displacement z from the equilibrium stretched
position y0 . Both y and z are positive in the figure.
Writing an ODE so that all occurrences of the dependent variable y are on the left hand
side and all other terms are on the right hand side, an ODE is said to be homogeneous if
the right hand side is zero and inhomogeneous if the right hand side is a constant or any
function of the independent variable, which in this case is t. Thus, Eq. (1.34) is seen to be
inhomogeneous.
The solution y(t) = a cos ω0 t + b sin ω0 t that worked for the homogeneous case does not
work here because of the inhomogeneous term. In general, solving inhomogeneous ODEs
is more difficult than solving homogeneous ODEs, but in this case a simple trick suffices to
solve the equation. To begin, we consider the case when the mass is simply at rest hanging
at its equilibrium position. In that case, there is no acceleration and Eq. (1.34) becomes
k
y=g. (1.35)
m
Solving for y, we obtain the equilibrium position y0 of the mass
mg
y0 = . (1.36)
k
We now define a new variable z = y − y0 = y − mg/k, which is the displacement of the
mass from its equilibrium position. This means that y = y0 + z = mg/k + z. Substituting
this expression for y into Eq. (1.34) yields
d2 z k
+ z=0, (1.37)
dt2 m
which is the same homogeneous equation of motion we obtained for the horizontal spring.
Therefore, the same general solution works: z(t) = a cos ω0 t+b sin ω0 t, which is equivalent
to y(t) = mg/k + a cos ω0 t + b sin ω0 t.
The trick for obtaining a homogeneous equation is to define the displacement z as the
displacement from the equilibrium position of the system, not from the position corre-
sponding to the unstretched spring. This illustrates a general principle for analyzing os-
cillations: always define your dynamic dependent variable as the displacement from the
equilibrium position. Because the system oscillates about its equilibrium position, this will
13 The ideal spring
t
t
Fig. 1.7 (a) A spring attached to two masses on an inclined plane showing l0 the equilibrium
length of the spring and x0 it equilibrium position with the two masses. (b) The same
spring after one mass has been removed.
always result in the simplest equation of motion. If you look back at our formulation of the
problem of oscillations of a diatomic molecule, you will see that we used this trick without
much fanfare to eliminate r0 from Eq. (1.26), which otherwise would have resulted in an
inhomogeneous equation.
Example 1.1 Two stacked blocks of mass m rest on a frictionless inclined plane as
shown in Fig. 1.7. The bottom mass is connected to the wall by a spring of unstretched
length l0 and spring constant k. Calculate the initial rest position x0 of the two stacked
masses relative to the wall.
At time t = 0 the top mass is suddenly removed. Calculate the position x(t) for times
t > 0 and determine the maximum value of x(t). The spring is stretched when x > l0 ; the
velocity is positive when the mass moves upward along the incline.
Solution
To find the initial rest position x0 of the two masses, we balance the forces in the x direction
along the inclined plane, noting that the spring force is directed down the plane for a
stretched spring when x − l0 > 0:
2mg
F x = −k(x − l0 ) − 2mg sin θ = 0 ⇒ x0 = l0 − sin θ.
k
Thus, at equilibrium, x0 < l0 ; the spring is compressed and exerts an upward force that
balances the downward gravitational force of the two masses.
Now at t = 0, the upper mass is removed. The equilibrium position is no longer x0 , and
the remaining block will move. Its equation of motion is
d2 x
m = −k(x − l0 ) − mgsinθ .
dt2
The initial conditions are x(0) = x0 and 3(0) = 0. We now follow the same procedure
outlined above for a vertical spring. First, calculate the new equilibrium position x1 by
setting the right hand side of the equation equal to zero: x1 = l0 − (mg/k) sin θ. Next, define
z(t) = x(t) − x1 . The equation of motion for z is now identical to Eq. (1.37) for a simple
harmonic oscillator with solution z(t) = a cos ω0 t + b sin ω0 t. The velocity at t = 0 is zero,
14 Oscillations
t
so b = 0, and z(0) = a = x0 − x1 = −(mg/k) sin θ. Putting this together gives
x(t) = x1 + z(t) = l0 − (mg/k) sin θ (1 + cos ω0 t).
The maximum displacement is thus x(T/2) = l0 .
1.3.3 Initial conditions and constraints
Most of the time, we use initial conditions to determine the two unknown constants in the
solution to the equation for a simple harmonic oscillator. However, these constants can be
determined in other ways. For example, position and velocity might be specified at some
time other than t = 0, or two positions might be specified at different times. Mathemati-
cally, we simply need two independent conditions about a system’s motion. Once they are
set, it is only a matter of algebra to determine the unknown constants in the general solution
and thus determine the motion at all times.
Usually the physical conditions of a system, often in the form of constraints that are
released at t = 0, determine the initial conditions. For example, in the vertical spring
example of §1.3.2, the physical condition was that the mass was held in place (constrained
to be) at the rest length of the spring for times t < 0 and then released at time t = 0 so that
y(0) = 0 and ẏ(0) = 0, implying z(0) = −mg/k and ż(0) = 0. In the example of the inclined-
spring, it is the compression of the spring by the two masses and the sudden removal of the
second mass that sets the constraint that x(0) = x0 and 3(0) = 0.
Collisions between a moving particle and a mass attached to a spring can lead to changes
in position and velocity and can thus set a system into oscillation. For example, consider a
particle of mass m moving with velocity 3 on a frictionless table; the mass m is moving to
the left initially so 3 < 0. At t = 0, it collides with a mass M at rest at the end of a spring
with force constant k, as shown in Fig. 1.8(a).
Consider first the case where the two particles stick together after they collide, as shown
in Fig. 1.8(b). The collision is completely inelastic, so energy is not conserved, but momen-
tum is. The velocity V 0 of the two particles at t = 0+ just after they collide is determined
by (m + M)V 0 = m3, which gives V 0 = m3/(m + M). The oscillator now has mass (m + M)
and natural frequency ω1 = [k/(m + M)]1/2 . The motion of the mass for t > 0 is
m3
x(t) = sin ω1 t , (1.38)
(m + M) ω1
where x = 0 is the initial and equilibrium position of mass M.

In an elastic collision, both energy and momentum are conserved. The two masses are
in contact just for an instant. If the initial velocities of masses M and m are, respectively,
V and 3, and their final velocities are V 0 and 30 , then
m3 + MV = m30 + MV 0
2
(1.39)
1
2 m3 + 12 MV 2 = 12 m302 + 12 MV 02 ,
15 Energy of a simple harmonic oscillator
t
(a) (b)
k
M m M m
(c) (d)
k V
M m M m
t
Fig. 1.8 (a) Mass m moving toward mass M before an inelastic collisions and (b) masses m
and M stuck together an moving with velocity after collision. (c) Masses M and m
approaching each other with respective velocities v and V before collision and (d)
masses m and M after the elastic collision.
which leads to
M−m 2m
V0 = V+ 3
M+m M+m (1.40)
2M m−M
30 = V+ 3.
M+m m+M
If V = 0, the initial condition is V 0 = [2m/(M + m)] 3. In this case the oscillator has a
mass of M, since the two masses separate after the collision, and the natural frequency is
√
ω2 = k/M.
As a final example, initial velocity may be determined by an impulse such as might be
delivered by a croquet mallet. An impulse is a force delivered over a short time period ∆t:
Z t+∆t
I= Fdt0 = ∆p (1.41)
t
where ∆p is the change in momentum brought about by the impulse. Thus the change in
velocity of a particle of mass m produced by an impulse is ∆3 = I/m.
1.4 Energy of a simple harmonic oscillator
Our analysis of the pendulum started by analyzing the force while our analysis of the
oxygen molecule started from the potential energy between the two atoms and from that
derived the force law. The two descriptions of oscillating systems are equivalent of course.
Here we take a closer look at simple harmonic oscillators from the perspective of the con-
servation of energy.
The change in potential energy ∆U associated with compressing or stretching a spring
16 Oscillations
t
is defined as the negative of the work W done by the spring in moving from some position
x1 to another position x2 :
Z x2
1
∆U = −W12 = − −kx dx = k(x22 − x12 ) . (1.42)
x1 2
Taking the zero of potential energy to correspond to the unstretched spring (x = 0), the
potential energy of a spring becomes
1 2
U(x) = kx . (1.43)
2
The principle of conservation of energy for the mass-spring system can be obtained by
applying Newton’s second law:
F = −kx = m ẍ , (1.44)
which can be rewritten as
d3
−kx = m . (1.45)
dt
Multiplying by dx and noting that dx = 3 dt, this can be written as
−kx dx = m3 d3 (1.46)
Integrating both sides gives
1 2 1 2
m3 + kx = E , (1.47)
2 2
where E is an integration constant. The two terms on the left hand side in Eq. (1.47) are
the kinetic and potential energies, respectively. Equation (1.47) states that their sum is a
constant E, which we identify as the total energy of the system.
We can also verify explicitly that energy is conserved for this system. The time-dependent
displacement x(t) of the mass on the spring is sinusoidal, just as for the pendulum and di-
atomic molecule,
x(t) = a cos(ω0 t − δ) , (1.48)
√
where ω0 = k/m. Writing 3 = ẋ and noting that k = mω20 , we calculate the total energy
at any given time t to be
1 2 1 2
E= m3 + kx (1.49)
2 2
1 1
= mω20 a2 sin2 (ω0 t − δ) + ka2 cos2 (ω0 t − δ) (1.50)
2 2
1
= mω20 a2 sin2 (ω0 t − δ) + cos2 (ω0 t − δ)
h i
(1.51)
2
1
= mω20 a2 . (1.52)
2
Thus, the total energy of the system is explicitly demonstrated to be independent of time.
Incidentally we remark that the total energy of any linear oscillator can be written in the
form of Eq. (1.52) as a mass times the oscillation frequency squared times the amplitude
squared.
17 Energy of a simple harmonic oscillator
t
E
K(x)
U(x)
x
−A A
t
Fig. 1.9 Kinetic energy K(x), potential energy U(x) and total energy E(x) = E as a function of
particle position.
Of course, this means that at each point along its trajectory, the total energy of the os-
cillator is constant. Figure 1.9 shows plots of the kinetic, potential, and total energy as a
function of the oscillator position x. The kinetic energy is a maximum at the equilibrium
position x = 0 while the potential energy is a maximum at the turning points x = ±a,
where the pendulum momentarily comes to rest and reverses direction. The plot illustrates
that the sum of the potential and kinetic energies is constant, consistent with Eq. (1.47).
Exercise 1.4.1 √Show that the absolute value of the position-dependent velocity is given
by |3| = ω0 a2 − x2 and that the kinetic energy K(x) of the mass-spring system is
given by K(x) = 12 mω20 (a2 − x2 ).
The potential energy U(x) of Eq. (1.43) is, of course, only the potential energy of the
stretched or compressed spring. The total potential energy will in general contain other
contributions. For example the vertical spring considered in §1.3.2 has an additional gravi-
tational potential energy Ug (y) = −mgy (with a negative sign because y is measured down-
ward) to yield a total potential energy of
1 2
UT (y) = ky − mgy, (1.53)
2
The potential energy here is at its minimum when dUT (y)/dy = 0, i.e.when y = y0 = mg/k.
An equilibrium state is one in which the force on the particle is zero or equivalently one
that is a minimum of the total potential energy.
Example 1.2 Figure 1.10 shows a block of mass m on a frictionless horizontal plane
attached to two identical springs with spring constants k and rest lengths l0 . The springs
are attached to fixed walls whose separation is such that in equilibrium, both springs are
stretched to a length 2l0 . What is the period of oscillation of the block: (a) if it is displaced
a distance x parallel to the springs and (b) if it is displaced a distance y perpendicular to
the plates?
Solution
There are two ways to solve this problem, and it worth looking at both methods. In the first
method, we calculate the potential energy U of springs as a function of x and y and then
18 Oscillations
t
t
Fig. 1.10 A block (which is interpreted as a point mass) on a frictionless table attached to two
stretched springs of rest length l0 with spring constant k.(a) Motion along the x
direction and (b) motion along the y direction.
calculate forces by differentiation. In the second, we calculate force directly, remembering

that the force is always along the direction of the stretched spring. The energy of the two
spring system is
1
U = k(|l1 − l0 |2 + |l2 − l0 |2 ),
2
where l1 and l2 are the lengths of springs 1 and 2, respectively.
(a) For motion along the x-direction, l1 = 2l0 + x and l2 = 2l0 − x. Then
1
U= k[(2l0 + x − l0 )2 + (2l0 − x − l0 )2 ] = k[l02 + x2 ],
2
√
F x = −dU/dx = −2kx = m ẍ, which gives an oscillation frequency of ω0 = 2k/m,
√
which corresponds to a period of T = 2π/ω0 = π 2m/k. Alternatively F x can be
calculated directly. The first string pulls in the direction that will tend to restore it to its
rest length l0 , i.e. in the minus-x direction. Its length is 2l0 + x (which is greater than l0
when x > 0), and it exerts a force F1 = −k(2l0 + x − l0 ). The second string pulls in the
positive-x direction to restore its rest length. The second string has length 2l0 − x, and
it exerts a force F2 = k(2l0 − x − l0 ). The total force from the springs is F1 + F2 = −2kx.
(b) For motion along the y-direction, l1 = l2 = (2l0 )2 + y2 .

p
s
!2
y
q
l1 = l2 = l = (2l0 )2 + y2 = 2l0 1 +
2l0
2
y y2
 
= 2l0 1 + 2 + · · ·  ≈ 2l0 + .
 
8l0 4l0
where we have used the Taylor series expansion discussed in §1.6 at the end of the
chapter to approximate the square root expression (see Eq. (1.89)). This is valid in the
limit of small amplitude oscillations where y l0 . Then
 !2 !4 
1 2 2 1 y 1 y  1
 , ≈ kl02 + ky2 ,

U(y) = k 2(l − l0 ) = kl0 1 + +

 (1.54)
2 2 l0 16 l0 2
where in the final step we dropped the small term proportional to (y/l0 )4 . Differenti-
√
ating gives the force: Fy = −dU/dy = −ky = mÿ and ω = k/m. Alternatively Fy
19 Complex exponential notation
t
can be calculated directly. Let T = k(l − l0 ) be the magnitude of the force that each of
the springs exerts on the displaced block. These forces are along the direction of the
displaced springs as shown in Fig. 1.10. The component of the combined force of the
two springs in the y direction is Fy = 2T sin θ ≈ −2T y/(2l0 ), where θ is the angle the
springs make with the x axis. The spring length l is 2l0 plus terms of order y2 . Thus to
obtain a force to linear order in y, we can replace l by 2l0 to yield T = kl0 and Fy = −ky
in agreement with the energy calculation.
If the equilibrium length of the springs were their unstretched length l0 rather than 2l0 ,
then the transverse displacements would be (l − l0 ) = y2 /(2l0 ) + · · · , and U(y) ≈ ky4 /(2l0 ).
This is an anharmonic spring, whose equation of motion is not so easily solved. This
example shows that it is the stretching of the springs, which puts them under tension, that
is responsible for the harmonic restoring force for transverse displacements. As we shall
see, this is essentially the same physics that controls the vibrations of a violin string.
1.5 Complex exponential notation
Simple harmonic oscillations are described by sinusoidal functions. Unfortunately, work-

ing with sine and cosine functions is mathematically cumbersome. However, there exists
an easier way to do algebra with sinusoidal functions, a way that involves the exponential
function with complex arguments. This section introduces you to this algebra and its use
in describing oscillations.
1.5.1 Complex exponentials
We introduce this approach√by considering the Taylor series expansion of exp(iα), where
α is a real number and i ≡ −1:
(iα)2 (iα)3 (iα)4 (iα)5
eiα = 1 + iα + + + + − ... (1.55)
2! 3! 4! 5!
α2 α3 α4 α5
= 1 + iα − −i + +i − ... . (1.56)
2! 3! 4! 5!
Collecting the individual terms in the sum in Eq. (1.56) into their real and imaginary parts,
which corresponds to collecting together the even and odd powers of α in the Taylor series,
gives
α2 α4 α3 α5
! !
iα
e = 1− + − ... + i α − + − ... . (1.57)
2! 4! 3! 5!
We recognize the series inside the two sets of parentheses as the Taylor series for cosine
and sine, respectively. Thus, Eq. (1.57) yields
eiα = cos α + i sin α . (1.58)
20 Oscillations
t
t
Fig. 1.11 The complex plane with complex number z represented in (a) cartesian coordinate
representation and (b) polar representation. Arrow shows the circle traced out when
α = ωt, in which case x = r cos ωt and y = r sin ωt.
Equation (1.58) is one of the most amazing and useful formulas in mathematics.1 It is also
the primary reason for introducing complex numbers into physics.
An alternative derivation of Eq. (1.58) is instructive. The complex exponential obeys the
same rules of differentiation that real exponentials do. Therefore,
d2 iα
e = i2 eiα = −eiα , (1.59)
dα2
and eiα obeys Eq. (1.5), the same differential equation as the displacement in a simple har-
monic oscillator, but with ω0 = 1. Thus eiα must be of the form of Eq. (1.6). To determine
α→0
the unknown coefficients a and b, we note that eiα |α=0 = 1 and deiα /dα = ieiα −−−→ i. Thus
a = 1 and b = i and Eq. (1.58) is reproduced.
Before delving into the uses of Eq. (1.58), let’s consider the graphical representation of
complex numbers. Any complex number z can be written in terms of its real and imaginary
parts: z = x + iy, where x and y are real numbers. Therefore any complex number can
be plotted in an (x, y) coordinate system known as the complex plane, as shown in Fig.
1.11(a). In this representation of complex numbers, the x coordinate gives the real part of
the complex number z and the y coordinate gives the imaginary part. The vector proceeding
from the origin to the point (x, y) in the complex plane is sometimes
p called a phasor.
We can also represent z in terms of its absolute value r ≡ |z| = x2 + y2 and the angle α
that it makes with the real (x) axis. Noting that x = r cos α and y = r sin α, we obtain
z = x + iy (1.60)
= r cos α + i r sin α (1.61)
iα
= re . (1.62)
Equation (1.62) is sometimes called the polar representation of z. Figure 1.11(b) shows
the polar representation of z in graphical form. We can convert between the cartesian and
polar representations of a complex number guided by the diagrams in Fig. 1.11. To convert
1 Setting α = π, Eq. (1.58) becomes, with a little manipulation, eiπ + 1 = 0, an economical and almost magical
relationship between the two most fundamental transcendental numbers, e and π, the multiplicative and
additive identities, 1 and 0, and the fundamental imaginary number, i.
21 Complex exponential notation
t
from the polar to the cartesian representation z = reiα , we use
x = r cos α (1.63)
y = r sin α . (1.64)
To convert from the Cartesian to the polar representation, we use
√ q
r = zz∗ = (x + iy)(x − iy) = x2 + y2
p
(1.65)
Im z y
!
α = tan−1 = tan−1 . (1.66)
Re z x
Figure 1.11 also provides us with a picture of how an oscillating sinusoidal function is
represented by the real part of a complex exponential. If we let the angle α = ωt, then as
t increases, z = reiωt sweeps out a circle of radius |z| = r moving counterclockwise in the
complex plane, as shown in Fig. 1.11(b), with a sinusoidal projection onto the real axis
given by x = Re(z) = r cos ωt.
1.5.2 Complex exponentials in harmonic oscillator problems
Armed with this knowledge about complex exponentials, how do we use them to describe
oscillations? There are several different options. Below, we introduce two that are widely
used and that we will find useful in our study of oscillations.
We start by pointing out that e−iωt is a solution to Eq. (1.5), the harmonic oscillator
equation, provided ω = ±ω0 , as you can readily verify. Thus, there are two independent
solutions, as required, and the general solution can be expressed by
θ(t) = a1 e−iω0 t + a2 eiω0 t . (1.67)
The coefficients a1 and a2 can in general be complex, which provide four independent
parameters. We can eliminate two of those parameters, by insisting that θ(t) be real, as we
would expect for physically meaningful solutions. Demanding that θ(t) be real means that
the imaginary part of θ(t) is zero. Thus, θ(t) and its complex conjugate are equal:
a1 e−iω0 t + a2 eiω0 t = a∗1 eiω0 t + a∗2 e−iω0 t . (1.68)
Collecting terms together gives
(a1 − a∗2 )e−iω0 t + (a2 − a∗1 )eiω0 t = 0 . (1.69)
Because e−iω0 t and eiω0 t are independent functions, the only way Eq. (1.69) can be zero is if
the coefficients (a1 − a∗2 ) and (a2 − a∗1 ) are zero. Therefore, requiring that θ(t) be real means
that a2 = a∗1 (or equivalently a1 = a∗2 ). This leaves only two independent parameters, the
amplitude and phase of a1 , a satisfying situation since we know that there are only two free
parameters in the general solutions. Expressing a1 = (A/2)eiδ in terms of its real amplitude
A/2 and its phase δ leads to
θ(t) = A cos(ω0 t − δ) , (1.70)
where we have used a2 = a∗1 = (A/2)e−iδ . Thus we see that Eq. (1.67) provides a completely
22 Oscillations
t
general solution equivalent to Eq. (1.10). The amplitude A and phase δ, or equivalently the
real and imaginary parts of a1 (= a∗2 ) are determined by the initial conditions.
Starting from the form given by Eq. (1.10) (or Eq. (1.70)), we can also write the general
solution to Eq. (1.5) as
θ(t) ≡ Re Ãe−iω0 t ,
h i
(1.71)
where we have introduced the complex amplitude Ã = Aeiδ . In this case, we use only one
of the two complex solutions, the real part, and thus have only two constants, the amplitude
A and phase δ of Ã.
We will find the complex form of Eq. (1.71) to be quite useful for the treatment of
cases when there is damping and when there are periodic external driving forces. There are
many other equivalent ways of expressing θ(t) in terms of a complex function. For example,
replacing ω0 and δ by −ω0 and −δ or replacing δ by δ − π/2 and taking the imaginary part
of the resulting expression provide representations that are equivalent to Eq. (1.71). We
will almost exclusively use Eq. (1.71). We will also find the complex form of Eq. (1.67) to
be useful in certain situations.
Exercise 1.5.1 The motion of a certain oscillator is described by the real part of the
equation
B
x(t) = eiωt (1.72)
−15 + 2i
where B is a real number. What are the amplitude and phase of x(t)? That is, rewriting
Eq. (1.72) in the form x(t) = A cos(ωt − δ), find expressions for the amplitude A and
phase φ, where A and φ are real numbers. Hint: find the magnitude r and the phase
α as illustrated in Fig. 1.11.
As another example, suppose you want to express cos(α+β) in terms of sines and cosines
of the angles α and β and not their sum. You could rack your memory, look this up, or you
could simply derive the correct expression using complex exponentials. Following the last
route, we first express cos(α + β) in terms of complex exponentials, remembering that we
are keeping only the real parts
cos(α + β) = ei(α+β) = eiα eiβ (1.73)
= (cos α + i sin α)(cos β + i sin β) (1.74)
= (cos α cos β − sin α sin β) + (1.75)
i (sin α cos β + cos α sin β) . (1.76)
Keeping only the real parts at the end of our calculation, we arrive at the purportedly
familiar formula
cos(α + β) = cos α cos β − sin α sin β . (1.77)
Using the complex exponential notation, there is no need to remember or look up such
formulas as they are readily derived in a few lines. In doing the calculation, however, it
is important to keep all the terms until the end to make sure that when you discard the
imaginary parts, you are not inadvertently discarding some real part of the expression.
23 Series expansions
t
Exercise 1.5.2 Following the same procedure shown above, find a similar formula for
sin(α + β) in terms of sines and cosines of the angles α and β. Hint: Explain why you
should start by writing sin(α + β) = −i ei(α+β) .
Exercise 1.5.3 Starting from the expression eiθ = cos θ + i sin θ, show that
1 iθ
sin θ = 2i (e − e )
−iθ
(1.78)
1 iθ
cos θ = 2 (e + e ) .
−iθ
(1.79)
Exercise 1.5.4 The hyperbolic sine and hyperbolic cosine functions are defined as
sinh x = 12 e x − e−x (1.80)

1 x
cosh x = 2 e + e −x
. (1.81)
Show that sinh(x) = −i sin(ix) and that cosh(x) = cos(ix).

Exercise 1.5.5 Consider the function
y(t) = Re a1 eiωt + a2 e−iωt ,

h i
(1.82)
where a1 and a2 are complex constants. Show that an alternative way of writing y(t)
is
y(t) = A1 cos ωt + A2 sin ωt , (1.83)
where A1 and A2 are real constants. In particular, derive expressions for A1 and A2
in terms of a1 and a2 . Using Hint: Write a1 = α1 − iβ1 and a2 = α2 + iβ2 where α1 ,
β1 , α2 , and β2 are all real. Expand the complex exponentials in Eq. (1.82) using Eq.
(1.58) and then find expressions for A1 and A2 in terms of α and β.
1.6 Series expansions
We have made extensive use of series expansions in this chapter, usually in order to sim-
plify the mathematical analysis of a problem. All of them can be derived from the general
expression for a Taylor series expansion of a function f (x):
f 00 (a) f 000 (a)
f (x) = f (a) + f 0 (a)(x − a) + (x − a)2 + (x − a)3 + ... (1.84)
2! 3!
Here we list series expansions for a few common functions. They can be derived from Eq.
(1.84) but as they are used so often, you should commit them to memory:
ex = 1 + x + 1 2 1 3 1 4
2! x + 3! x + 4! x + ... (1.85)
1 3 1 5
sin x = x − 3! x + 5! x − ... (1.86)
1 2 1 4
cos x = 1 − 2! x + 4! x − ... (1.87)
1 2 1 3 1 4
ln(1 + x) = x − 2 x + 3 x − 4 x + ... (1.88)
n
(1 + x) = 1 + nx + n(n−1) 2
2! x +
n(n−1)(n−2) 3
3! x + ... (1.89)
24 Oscillations
t
The Taylor expansion for (1 + x)n , known as the binomial expansion, is often overlooked
by new physics students. Nevertheless, it is one of the most useful expansions in physics. It
works for both positive and negative values of n, as well as for both integer and non-integer
values.
Example 1.3 Find the Taylor series expansion of f (y) = 1/ a2 − y2 .

p
Solution
We start by rewriting f (y) in the form (1 + y)n
1 1
f (y) = p = (1.90)
a2 y2 a 1 − (y/a)2
p
−
1
" y 2 #−1/2
= 1− . (1.91)
a a
Equation (1.91) is now in the form of Eq. (1.89) where x = −(y/a)2 and n = −1/2. Thus,
the Taylor series expansion is
 ! 2 ! ! 2 !2 
1  1 y 1 1 3 y
!
f (y) = 1 + − + + ...

− − − − (1.92)
a 2 a 2 2 2 a
1 1 y 2 3 y 4
" #
= 1+ + + ... (1.93)
a 2 a 8 a
Frequently, only the the first two terms are needed when using the binomial expansion.
Thus, it is handy to memorize it: (1 + x)n ' 1 + nx for |x| 1.
Exercise 1.6.1 Find the first two non-vanishing terms in the Taylor series for the fol-
lowing expressions (the first terms are given; you should fill in the “..." term). To do
so, first cast the relevant part of each expression in the form (1 + x)n , where |x| 1,
and then use Eq. (1.89) to perform the expansion. The expansion should be a poly-
nomial in the expansion variable similar to the example above.
1 Q
E x (x) = √ (1.94)
4π0 x x2 + a2
1 Q
' (1 + ...) where x a
4π0 x2
1
K(3) = (γ − 1)mc2 , where γ = p (1.95)
1 − 32 /c2
1
' m32 (1 + ...) v c
2
Equation (1.94) gives the electric field E x along the x axis generated by a uniformly
charged nonconducting wire of length 2a oriented along the y axis and carrying a
total charge Q. Equation (1.95) gives the relativistic expression for the kinetic energy.
25 Summary of important points of Chapter 1
t
Does each expression reduce to its expected result in the limits specified above?
Briefly explain (one sentence each!).
1.7 Summary of important points of Chapter 1
At the end of this and every chapter, we provide a summary of the chapter’s main points.
You should understand all of the points summarized below. The descriptions in the sum-
mary are kept brief deliberately. Refer back to the chapter for more extensive explanations
and for illustrative examples.
• Oscillations can occur when there is a minimum in the potential energy and the system is
displaced from that minimum. The restoring force, given by the gradient of the potential
energy, pushes the system towards equilibrium in the vicinity of the minimum.
• The restoring force is usually not linear, or equivalently, the potential energy is not a
parabola. A nonlinear force law leads to a nonlinear equation of motion, which in general
does not have an analytic solution (exceptions do exist, however). The first line of attack
for such a problem is to expand the force law in a Taylor series about the equilibrium
position and to keep only the linear term. The resulting equation of motion is linear and
has analytical solutions. The linearized force law can also be obtained by expanding the
potential energy in a Taylor series about the equilibrium position and keeping terms that
are quadratic in the displacement from the equilibrium position.
• Whenever you linearize a force law, you must take care to determine the range of dis-
placements over which the linearized force law provides an accurate description of the
system. An estimate for the displacement at which the error becomes bigger than the
desired accuracy (e.g. 1% or 10%) can be determined by keeping the next non-vanishing
term in the Taylor series. It is important to understand that some force laws can have a
relatively large range of displacements over which the linear force law is valid, while
other force law may have very small linear ranges. The linearized force law for pendu-
lum, for example, is accurate over a relatively large range of displacements, about one
third of a half revolution, while the linearized Morse potential for a diatomic molecule
is valid over displacements only on the order of 1% of the distance between atoms.
• The oscillatory motion that results from a linearized force law is called simple harmonic
motion.
• Obtaining solutions to problems involving oscillations about a linearized force law are
greatly simplified by using complex exponentials instead of sine and cosine functions.
• Because the linearized equation of motion describing free oscillations is a linear homo-
geneous second-order differential equation, it has two independent solutions, sine and
cosine, and the most general solution is a linear combination of these two solutions. The
general solution also has two integration constants, which, for a specific problem, can be
determined from the initial conditions. Specifying the initial conditions usually means
26 Oscillations
t
specifying the initial position and velocity of the system. If external forces, like a grav-
itational field, are added, the equation of motions becomes inhomogeneous. For simple
oscillators such as those treated here, the inhomogeneous equation can be converted into
a homogeneous equation by choosing the coordinates so that the equilibrium position of
the system corresponds to zero displacement.
27 Problems
t
Problems
1.1 A particle undergoing simple harmonic motion travels a total distance of 6.98 cm
in a full period T = 1.71 s. (a) What is the average speed of the particle? (ans:
4.08 cm/sec) (b) What are its maximum speed and acceleration? (ans: 6.41 cm/s,
23.6 cm/s2 )
1.2 A mass m = 0.2 kg attached to the end of a spring with spring constant k = 1 N/m is
released from rest at t = 0 s from an extended position xm . After 0.5 s, the speed of
the mass is measured to be 1.5 m/s. Calculate xm , the maximum speed, and the total
energy. (Ans: xm = 0.75 m, vm = 1.67 m/s, E = 0.28 J)
1.3 A volume V of air is contained in a vertical cylinder of radius R that is closed at the
bottom and fitted with piston of mass m at the top. The piston fits snugly so that no
air escapes the cylinder. Nevertheless, assume that the piston moves freely up and
down without friction. The equilibrium height of the bottom of the piston is h such
that the equilibrium volume of air is πR2 h.
(a) When the piston is displaced a distance x from its equilibrium position, show
that the restoring force is given by:
!γ
h
" #
F(x) = mg −1 (1.96)
h+x
where γ = C p /Cv = 1.4 for air. Assume that the piston is displaced from its
equilibrium position for a short enough time that no heat enters or leaves the
cylinder of gas.
(b) Expand the force law in a Taylor series about the piston’s equilibrium position
to quadratic order in the displacement x. What is the effective spring constant of
this system?
(c) Find an expression for the oscillation frequency of the piston for small displace-
ments about its equilibrium position. If the mass of the piston is 10 kg and the
equilibrium height is 0.10 m, what is the period of oscillation? Based on your
answer, do you think our assumption that there is no heat flow in or out of the
cylinder is justified?
1.4 In the text, we showed that the reduced mass µ = m/2 for a diatomic molecule made
of two identical atoms each of mass m. Show that the reduced mass of a diatomic
molecule made up of atoms of unequal masses m1 and m2 is given by µ = m1 m2 /(m1 +
m2 ). Verify that this reduces to the above result for the case that m1 = m2 . To do this,
follow the same methodology used in the text. Write down equations of motion for
each atom separately and then subtract one from the other to obtain an equation of
motion for the difference R2 − R1 , or for R2 − R1 − r0 . Hint: After obtaining the
equations for R1 and R2 separately, divide through by the appropriate mass before
subtracting the two equations to obtain an equation of motion for d2 (R2 − R1 )/dt2 .
28 Oscillations
t
1.5 The interaction between atoms in a molecule is sometimes approximated by the
Lennard-Jones2 potential
σ σ
" 12 6 #
U(r) = 4 − , (1.97)
r r
where σ and have units of length and energy, respectively, with values that depend
on the particular atoms involved.
(a) Plot the Lennard-Jones potential as a function of r. More specifically, plot U(r)/
as a function of r/σ for the range 0.9 ≤ r/σ ≤ 3. Scale the vertical axis so the
energy well is clearly visible.
(b) Show that the minimum energy for the Lennard-Jones potential is located at r =
r0 ≡ 21/6 σ. Physically, what does this distance represent? What is the minimum
energy? What is the binding energy for this potential?
(c) Show that the spring constant for the Lennard-Jones potential is k ' 57.1/σ2
and the oscillation frequency by ω ' 7.56 /µ/σ where µ is the reduced mass
p
for the system.

(d) On the same plot that you plotted the Lennard-Jones potential, plot its harmonic
approximation, which includes all the terms in the Taylor expansion of U(r)
about r = r0 up to and including the term quadratic in r − r0 .
(e) The vibrational frequency f of a nitrogen molecule is about 8.2 × 1013 Hz. Using
numbers found on the web or in books, determine the values of and µ. Using
those values, use the results from part (c) to find the the width of the well σ?
What is the value of the ratio σ/r0 ?
1.6 (*) A particle of mass m moves in a potential
U(x) = V0 (1 − cos αx) − mgx.
(a) What is its equilibrium position x0 (nearest to x = 0)? Express you answer in
terms of V0 , α and mg. You may assume that mg < αV0 . (ans: x0 = (1/α) sin−1 (mg/αV0 ))
(b) What is the frequency of small oscillations about the equilibrium position. Your
answer may be expressed as a function of A, α, m, and x0 . (ans: ω = α[(V0 /m) cos αx0 ]1/2 =
(α/m)1/2 [(α2 V02 − m2 g2 ]1/4 )
1.7 A block of mass m = 0.1 kg attached to an ideal spring moves on a frictionless
surface. Let x(t) be its displacement as a function of time.
(a) At time t = 0, the block passes through the point x = 0 moving to the right (i.e.,
toward positive x). At time t = 0.5 s, the block reaches its maximum excursion
of xm = 10 cm.
(i) What is the period T of oscillation? (Ans: t = 2 s)
(ii) What is the spring constant k? (Ans: k = 0.1π2 N/m)
(iii) What is the maximum velocity 3m ? (Ans: 0.314 m/s)
(iv) Write the full expression for x(t).
2 Sir John Lennard-Jones (1894–1954) was a professor of theoretical chemistry at Cambridge University. He
was one of the pioneers of molecular orbital theory.
29 Problems
t
t
Fig. 1.12 (a) Block at rest prior to time t = 0. (b) Configuration of the block at t = 0+ (i.e., right
after t = 0) shown in lighter gray and configuration for general time t in darker gray.
√
(b) At time t1 = 5T/8, a particle of mass m and moving with velocity 3 = 33m / 2 to
the right collides completely inelastically with the block (i.e., collides and sticks
to the block) of the oscillator in (a).
(i) What is the position of the block
√ as a function of time after the
√ collision?
√ 2π
(Ans: x(t) = − (3/2)xm cos[ T (t − t1 ) + φ] where φ = tan −1
2).
(ii) At what time does the block pass the origin again? (Ans: 1.53 s).
1.8 A block of mass m, resting in its equilibrium position on a frictionless platform, is
attached to a Hooke’s-law spring with spring constant k with one end fixed on a
vertical wall connected to the platform as shown in Fig. 1.12. At time t = 0, the
platform is suddenly tilted upward to an angle θ.
(a) Determine the position x(t) of the block for all t > 0. (ans: x(t) = l0 −(mg/k) sin θ (1−
√
cos ω0 t), where ω0 = k/m)
(b) What are the total kinetic energy K(t) and potential energy U(t) for times t > 0.
At what times is the kinetic energy a maximum? (ans: t = (1 + 2n)(T/4), where
T is the period of oscillation)
1.9 A cylindrical cork of radius R = 1.5 cm, height d = 0.75 cm, and specific gravity of
0.24 floats at an air-water interface.
(a) How far below the surface of the water is the bottom of the cork in equilibrium:
(ans: h = 0.18 cm).
(b) At time t = 0, a small of aluminum bead of mass mA = 1.0 gm is placed on
top of the cork while imparting a downward velocity of 30 = 10 cm/s to the
combined cork and bead system. Neglecting any friction or water flow induced
by the motion of the cork, calculate the depth y(t) of the bottom of the cork for
all t > 0 and sketch its form. What is the maximum depth and at what time is
it first reached? (ans: 0.55 cm at time 0.04s) What is the minimum depth and at
what time is it first reached? (ans: 0.09 cm at time 0.098s).
1.10 (*) Four identical springs of rest length l0 are attached to a mass m at right angles
to each other on a frictionless horizontal plane as shown below. The two springs
parallel to the x-axis are stretched to length 2l0 while the two springs parallel to the
y-axis have their rest length l0 . When the mass is attached to any one of the individual
springs alone, it oscillates with angular frequency ω0 = 2π/T 0 . Express your answers
in terms of ω0
30 Oscillations
t
t
Fig. 1.13 Figure for problem 1.10
(a) Calculate
√ the frequency ω x of small oscillations in the x direction. (Ans: ω =
2ω0 )
Calculate the frequency ωy of small oscillations in the y direction. (Ans: ω =
(b) √
3ω0 )
(c) At time t = 0, the block is given velocities 3 x = 30 and 3y = 30 in the x and y
directions, respectively. What are x(t) and y(t) for t > 0.
1.11 A mass m is attached to two walls a distance l0 away by massless linear springs
whose unstretched length is l0 as shown below in gray. The force of gravity pulls the
mass down and stretches the springs as shown below in black.
t
Fig. 1.14 Figure for problem 1.11
(a) Determine the equation for the equilibrium position y0 of the mass. (Ans: y0 is
the solution to the equation, T l(y) sin θ = k(l(y) − l0 )y, where l(y) = (l02 + y2 )1/2 ).)
(b) What is the frequency of oscillation of the mass about its equilibrium position?
Your answer can be expressed in terms of l(y0 ). (Ans: ω2 = k(y0 )/m, where
k(y0 ) = k[l3 (y0 ) − l03 ]/l3 (y0 ).)
(c) The same arrangement of springs and masses is placed on a horizontal friction-
less plane. Gravity does not operate. What is the potential energy of displacement
to lowest order in y? What is the equation of motion of the mass to lowest order
in y? (Ans: mÿ = −2ky3 /l02 )
2 Damped Oscillations
Real mechanical oscillations, like those of the pendulum or mass on a spring, die out with
time. This is because there is always some source of friction or damping in macroscopic
systems that robs an oscillator of energy. The damping of oscillations can be undesirable,
as for the case of a pendulum in a clock, or they can be purposely built into the design of a
system, as for a swinging door or the suspension of an automobile. In the case of a clock’s
pendulum, the aim is relatively simple: minimize damping. In other cases, however, the
damping needs to be set just right. The swinging door needs to swing, but its oscillations
should die out quickly. If the damping is too weak, the swinging of the door can surprise,
and harm, the next person coming through. If the damping is too strong, the door may
take too long to close or be too hard to open. Studying the physics of damped oscillations
helps us design oscillators that have the right degree of damping required for any particular
application.
2.1 The damped oscillator
We begin our discussion of damped oscillations with a mass suspended from the ideal
spring we studied in Chapter 1, but this time we include the effect of friction, which slows
the mass and causes its oscillations to die away. Friction can arise from air resistance.
Alternatively, we could suspend the mass in a liquid to provide more substantial damping.
Experiments reveal that these sources of damping can usually be accounted for quite well
by a damping force Fd that is proportional to the instantaneous velocity 3 of the mass
Fd = −Γ3 , (2.1)
where Γ is an experimentally-determined system-dependent damping or friction constant.

The minus sign indicates that the damping force acts in the direction opposite to the veloc-
ity of the mass (Γ > 0), so its effect is always to slow the motion of the pendulum.
We obtain the equation of motion for the damped harmonic oscillator following the same
procedure we used in Chapter 1, but now including the frictional force given by Eq. (2.1),
d2 y dy
m 2
= −ky − mg − Γ . (2.2)
dt dt
Following the procedure discussed in §1.3.2, we define z(t) = y(t)−mg/k to remove gravity
31
t
t
Fig. 2.1 Mass suspended from a spring and immersed in a fluid that damps its motion. As
shown the mass is stretched a distance z from its equilibrium position of y0 = mg/k.
The damping force Fd = −Γ3 is in the opposite direction of the velocity 3.
from the equation of motion, to obtain

d2 z dz
m = −kz − Γ . (2.3)
dt2 dt
Dividing through by m, we obtain the linearized equation of motion for the damped har-
monic oscillator:
d2 z dz
2
+γ + ω20 z = 0 , (2.4)
dt dt
√
where ω0 ≡ k/m, and γ = Γ/m. For the case of no damping where γ = 0, solutions to
Eq. (2.4) are sinusoidal. More precisely, the solutions are given by z = A cos(ω0 t − δ), so
we see that ω0 is just the oscillation frequency of the undamped oscillator. The amplitude
A and the phase δ are constants determined by the initial conditions.
When the damping constant γ is small, we expect the solutions still to be approximately
sinusoidal but with an amplitude that gets smaller with the passage of time. Based on these
physical considerations, a plausible guess for the solution would be a function of the form
z(t) ∼ A e−αt cos(ωt − δ) . (2.5)
Equation (2.5) is an oscillating solution with an amplitude a exp(−αt) that decays with
time. The constant α must have units of inverse time so that the argument αt of the expo-
nential is dimensionless. We might guess that α is proportional to γ because γ has units of
inverse time and, more importantly, because increasing γ increases the damping force (see
Eq. (2.1)), which should cause the oscillations to decay more rapidly.
Our guess, based on physical grounds, of Eq. (2.5) as a solution for the damped pen-
dulum turns out to be correct, as you can readily check by substituting Eq. (2.5) into Eq.
(2.4). Doing so yields explicit expressions for α and ω: α = γ/2 and
q q
ω ≡ ω1 = ω20 − (γ/2)2 = ω0 1 − (γ/2ω0 )2 , (2.6)
which is correct in the limit of weak damping, i.e. when 1 − (γ/2ω0 )2 > 0. However, the
solution fails in the limit of strong damping, i.e. when 1 − (γ/2ω0 )2 < 0. This approach can
33 The damped oscillator
t
be salvaged, but the method is mathematically cumbersome, primarily because it involves
sine and cosine functions.
We need a more systematic procedure for finding solutions to Eq. (2.4). While different
approaches are available, we shall choose one that is useful not just for the problem at
hand, but one that finds much wider application in other problems, specifically when we
consider more complex oscillators and waves. As we expect the solutions to be oscillatory,
we look for solutions of the form
z(t) = a e−iωt (2.7)
where a and ω can be complex numbers. By allowing ω to be complex, Eq. (2.7) can
express both sinusoidally oscillating and exponentially damped solutions, depending on
the value of ω. For example, suppose we let ω = b − ic, where t, b, and c are real positive
numbers. Then,
e−iωt = e−i(b−ic)t = e−ct e−ibt = e−ct (cos bt − i sin bt) , (2.8)
where we have used Eq. (1.58) with α = −bt to write
e−ibt = cos bt − i sin bt . (2.9)
Equation (2.8) is a decaying sinusoidal function, as promised. It’s true that the function
is complex, but this turns out not to be a problem. As discussed in §1.5, we shall use the
complex functions of the form eiα in order to simplify the analysis, but in the end will keep
only the real part of the solution. The imaginary part comes along for the ride, but is not
used.
The amplitude a in Eq. (2.7) can also be complex, which allows us to adjust not only
the real amplitude of z(t) but also the phase of its oscillations. To see this, we write the
complex amplitude a in polar form as discussed in §1.5
a = Aeiδ , (2.10)
where A and δ are real numbers. Letting ω = b − ic as before, the expression for z(t) in Eq.
(2.7) becomes
z(t) = a e−iωt = Aeiδ e−ct e−ibt = Ae−ct e−i(bt−φ)
= Ae−ct [cos(bt − φ) − i sin(bt − φ)] (2.11)
Taking only the real part of Eq. (2.11), z(t) = Ae−ct cos(bt − φ), we see that writing Eq.
(2.7) with a complex amplitude a and complex frequency ω allows us complete freedom to
express an exponentially decaying sinusoidal function with arbitrary (real) amplitude and
phase.
Armed with Eq. (2.7), we return to our task of solving the equation of motion for the
damped harmonic oscillator. To this end, we propose Eq. (2.7) as a general solution to
Eq. (2.4). Our task is to determine which sets of values of the constants a and ω in Eq.
(2.7) yield meaningful solutions to Eq. (2.4). Writing the time derivatives as ż ≡ dz/dt and
z̈ ≡ d2 x/dt2 , Eq. (2.4), the linearized equation of motion, becomes
z̈ + γ ż + ω20 z = 0 , (2.12)
t
The form of our proposed solution is z(t) = a exp(−iωt), which gives ż = −iωa exp(−iωt)
and z̈ = −ω2 a exp(−iωt). Substituting these expressions into Eq. (2.12) and then canceling
the common factors of a exp(−iωt) yields the equation
−ω2 − iγω + ω20 = 0 . (2.13)
The values of ω satisfying Eq. (2.13) are obtained from the quadratic formula
r
γ γ 2
ω± = −i ± ω20 − . (2.14)
2 2
This result leads to three different solutions depending on whether the value of the dis-
criminant ω20 − (γ/2)2 is positive, negative, or zero. These correspond to different degrees
of damping: underdamped, overdamped, and critically damped. We treat each of the three
cases in the sections that follow.
2.1.1 Underdamped response: (γ/2)2 < ω20
We begin our analysis with the case when (γ/2)2 < ω20 , that is when the rate of damping
is smaller than the rate of oscillations. In this limit, you might expect the oscillations to
proceed more or less unscathed, but to slowly die out over time. As we shall see, this is
indeed the case, at least in the extreme underdamped limit when (γ/2)2 ω20 .
When (γ/2)2 < ω20 , we can rewrite Eq. (2.14) as
ω± = −iγ/2 ± ωu , (2.15)
where ωu is a positive real frequency given by

q
ωu = ω20 − (γ/2)2 > 0 . (2.16)
Equation (2.15) actually implies that there are two independent solutions, one for the
plus sign and another for the minus sign in Eq. (2.15):
zu1 (t) = a1 e−iω+ t = a1 e−γt/2 eiωu t (2.17)

−iω− t −γt/2 −iωu t
zu2 (t) = a2 e = a2 e e (2.18)
Aside from the amplitudes a1 and a2 , the two solutions differ only by a minus sign in the
complex exponential, but that is enough to make the two solutions independent, meaning
that zu1 (t) , Czu2 (t) for any (complex) constant constant C. Note also that zu1 (t) , zu2 (−t).
Our next step is to exploit one of the most important properties of the solutions to linear
differential equations, namely, the principle of superposition. The principle of superposi-
tion states that if there are two (or more) solutions to a linear differential equation, then
any linear combination, i.e. any sum of two or more solutions, is also a solution to the
differential equation in question. Thus, because Eqs. (2.17) and (2.18) are both solutions
to the linearized equation of motion, Eq. (2.4), so is their sum. Therefore, the most general
t
solution to Eq. (2.4) for this problem is
1 1
zu (t) = a1 e−γt/2 eiωu t + a2 e−γt/2 e−iωu t (2.19)
2 2
1 −γt/2
a1 e + a2 e−iωu t .
iωu t

= e (2.20)
2
This is the general solution with coefficients a1 and a2 with four independent parameters
that are still undetermined. As in the case of complex solutions for the undamped oscillator
discussed in §1.5.2, requiring that our solution zu (t) be real means that a1 = a∗2 , which
eliminates two of the undetermined parameters and leaves only two, consistent with our
expectations for a second order differential equation. Writing a1 = a∗2 = Ae−iδ , where A
and δ are real, and substituting into Eq. (2.20) gives
1 −γt/2 i(ωu t−δ)

+ e−i(ωu t−δ)

zu (t) = Ae e (2.21)
2
= Ae−γt/2 cos(ωu t − δ) (2.22)
This equation can, of course, be expressed in terms of sine and cosine functions:
zu (t) = e−γt/2 (A1 cos ωu t + A2 sin ωu t) , (2.23)
where A1 = A cos δ and A2 = A sin δ. Equation (2.22) or alternatively Eq. (2.23) express
the general solution to the damped oscillator.
As in the case of the undamped oscillator, the initial conditions determine the two un-
known coefficients (A and δ or A1 and A2 ). As a specific example, let’s consider the case
where the oscillator is initially at rest, but has an initial velocity 30 . We apply these initial
conditions to the general solution given by Eq. (2.23)
zu (0) = A1 = 0 (2.24)
żu (0) = 30 = ωu A2 , (2.25)
which gives
30 −γt/2
zu (t) = e sin ωu t.
ωu
3i e−γt/2 q
= q sin ω20 − (γ/2)2 t . (2.26)
ω20 − (γ/2)2
The solution given by Eq. (2.26) is a damped sine wave with a frequency ωu given by Eq.
(2.16). It is plotted in Fig. 2.2. Note than ωu is smaller than the natural frequency ω0 of the
undamped pendulum. This result is consistent with what you might have expected at the
outset, namely that weak damping causes the oscillations to die off with time and also to
slow down somewhat.
There are four (at least!) equivalent ways that the general solution for the underdamped
t
z(t)
1.0 Q = 10
Q=3
Q=1
0.5
0.0 t/T
2 4 6 8
-0.5
-1.0
t
Fig. 2.2 Transient response of an underdamped oscillator for different values of Q. The gray
line shows the decaying exponential envelope of the Q = 10 trajectory.
oscillator can be written. We summarize them here:

iωu t ∗ −iωu t
 −γt/2



 e a 1 e + a 1 e

e−γt/2 (A1 cos ωu t + A2 sin ωu t)



zu (t) =  . (2.27)

B e−γt/2 cos(ωu t − δ)






C e−γt/2 sin(ω t − φ)


u
The solutions correspond to different ways of expressing a sinusoidal function with an

amplitude that decays exponentially in time. Each function has two constants that are de-
termined by the initial conditions for any given problem. In each case, the effect of the
constants is simply to adjust the phase of the sinusoidal oscillations within the decaying
exponential envelope. To solve a problem with any given set of initial conditions, you are
free to choose any of the above functional forms.
The quality factor of an oscillator

Oscillators are often characterized by a quantity known as the quality factor, defined here
as
ω0
Q= . (2.28)
γ
For the moment, we are considering the case when (γ/2)2 < ω20 , which means that Q > 21 .
The quality factor Q also roughly corresponds to the characteristic number of oscillations
a free oscillator undergoes before the oscillations effectively die out. Therefore, if a pen-
dulum oscillates 20 times back and forth before the amplitude of its oscillations diminish
appreciably, then Q for that pendulum is about 20 (to within a factor of 2 or so). Figure 2.2
shows the trajectories of a damped oscillator for several different values of Q.
t
From Eq. (2.16) it is also easy to show that the oscillation frequency for an underdamped
pendulum is given by ωu = ω0 [1 − (2Q)−2 ]1/2 . In the limit that Q 1, this means that the
change in the oscillation frequency from the undamped value is ∆ω = ωu −ω0 ' −(8Q2 )−1 ,
which is very small indeed when Q 0. Thus, a weakly damped pendulum oscillates at
very nearly the natural frequency.
When (γ/2)2 < ω20 , the system is said to be underdamped. This means simply that the
system will oscillate back and forth at least once, and usually many more times. For the
opposite case, when (γ/2)2 > ω20 , the system does not oscillate but simply relaxes directly
to its equilibrium position. In that case, the system is said to be overdamped. We examine
that case in the next section.
Exercise 2.1.1 For the underdamped case, show that the amplitude of the oscillations is
damped by about a factor of 20 after Q oscillations (i.e. for t = QT , where T = 2π/ω0
is the period).
2.1.2 Overdamped response: (γ/2)2 > ω20
For the overdamped case, the discriminant ω20 − (γ/2)2 in Eq. (2.14) is negative, which
means that the frequency ω is purely imaginary. To make this explicit, we rewrite Eq.
(2.14) as
 r 
 γ γ 2
2
ω± = −i  ± − ω0 

(2.29)
2 2
= −i (γ/2 ± β) , (2.30)
where
q
β= (γ/2)2 − ω20 > 0 . (2.31)
As in the underdamped case, Eq. (2.30) implies that there are two solutions. Once again,
the most general solution is a linear combination of the two solutions:
1 1
zo (t) = b1 e−( 2 γ+β)t + b2 e−( 2 γ−β)t (2.32)
= e−γt/2 b1 e−βt + b2 eβt ,

(2.33)
where b1 and b2 are constants that are determined by the initial conditions. Note that in this
case, there is no need for complex coefficients b1 and b2 .
Once again, as an example, let’s work out the solution for the initial conditions where
the oscillator starts at z = 0 with an angular velocity ż(0) = 30 . According to Eq. (2.32),
the initial position and velocity are given by
zo (t = 0) = b1 + b2 = 0 (2.34)
żo (t = 0) = −( 12 γ + β)b1 − ( 12 γ − β)b2 = 30 . (2.35)
Solving Eqs. (2.34) and (2.35) simultaneously gives b1 = −b2 = −Ωi /2β. Substituting this
t
z(t)
0.4
Q = 0.5
Q = 0.3
0.2
Q = 0.1
0.0 t/T
1 2 3 4
t
Fig. 2.3 Transient response of an overdamped oscillator for different values of Q. Note the
difference in time scale and the smaller-amplitude response compared to the
underdamped cases shown in Fig. 2.2. The critically damped case, which
corresponds to Q = 1/2, is the most rapidly decaying solution.
into Eq. (2.32) gives1

30 −γt/2 βt
zo (t) = e e − e−βt (2.36)
2β
30
= e−γt/2 sinh βt (2.37)
β
!
3i
q
= q e−γt/2 sinh (γ/2)2 − ω20 t (2.38)
(γ/2)2 − ω20
Note that this solution never crosses zero; it merely rises from zero to some maximum and
then decays back down to zero. Such a system is said to be overdamped. The condition for
the overdamped case is (γ/2)2 > ω20 , which corresponds to Q < 21 . Figure 2.3 shows the
trajectories of overdamped as well as underdamped oscillators. Note that while increasing
the amount of damping (by increasing γ) leads to a smaller maximum displacement for the
different cases in Fig. 2.3, it also has the effect of slowing the approach to equilibrium. The
critically damped pendulum reaches equilibrium the fastest.
Exercise 2.1.2 You can show that the general solutions of the overdamped and under-
damped cases are equivalent. Begin by noting that ωu = iβ. Then show that setting
ωu in Eq. (2.26) equal to iβ yields Eq. (2.37). Thus, the solutions to the underdamped
and overdamped cases are obtained from the a single set of equations.
Exercise 2.1.3 For t 2/γ, show that in the strongly overdamped limit, meaning that
1 Here we use the hyperbolic sine function, denoted sinh, and defined as sinh x ≡ 21 (e x − e−x ). There are
corresponding hyperbolic cosine and tangent functions defined as cosh x ≡ 12 (e x + e−x ) and
tanh x = sinh x/ cosh x.
t
Q 1, Eq. (2.38) can be written as
30 −γQ2 t
zo (t) ' e . (2.39)
γ
Describe in words how increasing the damping constant by a factor of 2 changes
the response of the of a damped oscillator in the strongly overdamped limit. Hint:
Keep only the positive exponential in the sinh function (why is this justified?) and
then expand the argument of the exponential using the binomial expansion (see Eq.
(1.89)).
The analysis above suggests two equivalent ways to write the general solution to the
over damped case, which we summarize them here for reference:
b1 eβt + b2 e−βt
 −γt/2
e

zo (t) =  . (2.40)

e−γt/2 (B cosh βt + B sinh βt)

1 2
Writing out the sinh βt and cosh βt in their equivalent exponential forms, you can readily
show that B1 = 21 (b1 + b2 ) and B2 = 12 (b1 − b2 ). You are free to use either form when
solving problems.
2.1.3 Critically damped oscillations: (γ/2)2 = ω20
For the critically damped case, the discriminant ω20 − (γ/2)2 = 0 and ω in Eq. (2.14) is
given by
γ
ω = −i (2.41)
2
From the form of our proposed solution given in Eq. (2.7), one may be temped to write the
solutions as
zc (t) = c e−γt/2 , (2.42)
and be done with it. But there is a problem with Eq. (2.42): It cannot describe the solution to
the situation where the initial displacement zc (t = 0) is zero and the initial angular velocity
żc (t = 0) is finite (i.e. not zero). Another clue that there is a problem comes from the fact
that only one integration constant, C, appears in Eq. (2.42). As a solution to a second order
differential equation, we expect two integration constants.
So how do we obtain the correct solution for the critically damped case? There are
several approaches available but perhaps the most obvious is to use either the underdamped
or overdamped solutions we have already found, and to take the limit of those solutions
as the discriminant (γ/2)2 − ω20 goes to zero. After all, if we make only an infinitesimal
change in the relative values of γ and ω0 , we expect only an infinitesimal change in the
solution, so such a procedure should lead to a well-defined result.
In fact, this procedure works very well. Problem 2.2 guides you through the mathemat-
ics. The result is
zc (t) = (c1 + c2 t) e−γt/2 , (2.43)
t
z(t)
1.0
critically damped
0.5
overdamped
underdamped
0.0 t/T
1 2 3 4
-0.5
t
Fig. 2.4 Three cases of damped oscillations: underdamped (solid), overdamped (short
dashes), and critically damped (long dashes).
where c1 and c2 are the two (expected) integration constants whose values are fixed by the
initial conditions.
Once again, as an example, let’s work out the solution for the initial conditions where
the oscillator is initially at z(0) = 0 with a velocity of ż(0) = 30 . In this case, the solution
is quite straightforward: z(0) = 0 implies that c1 = 0 in Eq. (2.43), while ż(0) = 30 implies
that c2 = 30 . Thus, the specific solution for these initial conditions is
zc (t) = 30 t e−γt/2 . (2.44)
Starting at t = 0 with an initial velocity of Ωi , the oscillator moves away from its equilib-
rium position, according to Eq. (2.44), and reaches a peak displacement time t = 2/γ, after
which the displacement decays exponentially towards zero.
2.1.4 Summary of damped oscillations
In Fig. 2.4, we summarize the three cases of damping treated above: underdamped, over-
damped, and critically damped. Here we show the case where each oscillator is started
at rest from a finite initial displacement. While derived in the context of the the damped
oscillator, these results are completely general and apply to all linear simple harmonic
oscillators. Notice that only the underdamped case exhibits oscillations. Figure 2.4 also
shows that an overdamped oscillator initially decays more rapidly than a critically damped
oscillator, but ultimately decays more slowly. This is an important consideration when de-
signing dampers for door as well as many other practical devices where oscillations are
undesirable. For these initial conditions, adjusting the damping constant γ so that it cor-
responds precisely to the case of critical damping ensures that the system returns to the
equilibrium position the fastest.
41 Energy loss from damping
t
2.2 Energy loss from damping
The amplitude of the oscillations of a damped oscillator decreases with each successive
swing as energy is lost to the friction of the damping force. The energy lost is simply the
work done by the damping force on the oscillator
Z
W= −Fd ds , (2.45)
where the minus sign arises according to Eq. (2.1) Fd = −Γ3 = −γm3. Substituting this
into Eq. (2.45) and using the fact that ds = 3 dt gives the work W(t) done by the damping
force since time t = 0
Z t
W(t) = −γm 32 (t0 ) dt0 . (2.46)
0
The work-energy theorem tells us that the change in energy is the work done on the system
by outside forces, in this case the friction force. Thus
Z t
E(t) − E0 = W(t) = −γm 32 (t0 ) dt0 . (2.47)
0
The instantaneous rate of energy loss is obtained by differentiating Eq. (2.47)

dE
= −γm 32 (t) . (2.48)
dt
Thus, the instantaneous rate of energy loss is proportional to 32 .
Of course the total energy can also be calculated directly from the instantaneous sum of
kinetic and potential energy,
1 2 1
E(t) = K(t) + U(t) = m3 (t) + kx2 (t) . (2.49)
2 2
Calculating E(t) is fairly tedious for the most general case, but simplifies considerably in
the limit of very weak or very strong damping, that is, in the strongly underdamped and
overdamped regimes. It also becomes a bit easier to calculate the results if we consider
specific initial conditions. So we use the initial conditions introduced previously of zero
initial displacement and a finite initial velocity. In the very weak damping limit, which
corresponds to (γ/2)2 ω20 or Q 1, we can use Eq. (2.26) and set ωu ' ω0 . This gives
30 −γt/2
z(t) = e sin ω0 t. (2.50)
ω0
Differentiating this expression, we obtain the velocity
30 −γt/2 h
ω0 cos ω0 t − 12 γ sin ω0 t .
i
3(t) = lż(t) ' e (2.51)
ω0
Because γ/2 ω0 , we can neglect the sine term above. Thus
3(t) = ż(t) ' 30 e−γt/2 cos ω0 t . (2.52)

t
Substituting Eqs. (2.51) and (2.52) into Eq. (2.49) gives
1 2 1
E(t) = m3 (t) + kx2 (t) (2.53)
2 2
1 2 −γt 1
' m30 e cos2 ω0 t + m320 e−γt sin2 ω0 t (2.54)
2 2
1 2 −γt
= m30 e (2.55)
2
1 2 −ω0 t/Q
= m30 e . (2.56)
2
Thus we see that in the very weak damping limit, energy decays exponentially in time with
a time constant of 1/γ. Or equivalently, since Q = ω0 /γ, energy decays on a scale that is
Q/2π times slower than the period of an oscillation.
In the strong damping limit, which corresponds to (γ/2)2 ω20 or Q 1, we start with
Eq. (2.37)
30 −γt/2
z(t) = = e sinh βt (2.57)
β
30 −γt/2 βt
= e e − e−βt (2.58)
2β
30 h (β−γ/2)t i
= e − e−(β+γ/2)t . (2.59)
2β
Except for an initial short transient, exp(−(β + γ/2)t) exp((β − γ/2)t). Therefore, to
determine the long-time behavior, we need to keep only the first exponential in Eq. (2.59).
We can further simplify matters by expanding the quantity (β − γ/2), keeping in mind that
(γ/2)2 ω20 :
γ γ
q
β− = (γ/2)2 − ω20 − (2.60)
2 2
s 
2
γ 
 4ω0 
=  1 − 2 − 1

(2.61)
2 γ
γ  1 4ω20
  
= 1 − + ... − 1 (2.62)
 
2 2 γ 2
ω20
'− (2.63)
γ
where in going from Eq. (2.61) to Eq. (2.62) we have used the binomial expansion (1−x)a =
1 − ax + O(x2 ). With these approximations, Eq. (2.59) becomes
30 −(ω2 /γ)t
z(t) ' e 0 . (2.64)
2β
Differentiating, we obtain
30 ω20 −(ω2 /γ)t
3(t) ' − e 0 . (2.65)
2β γ
43 Shock absorbers
t
Substituting these two expressions into Eq. (2.49) yields
1 2 1
E(t) = m3 (t) + kx2 (t) (2.66)
2 2
1 2 2 1 2 2
' m30 Q ( 4 Q + 1) e−2(ω0 /γ)t (2.67)
2
1 2 2 −2Qω0 t
' m30 Q e . (2.68)
2
Thus, in the strong damping limit where Q 1, the energy decays at a rate of 2Qω0 while
in the weak damping limit where Q 1, it decays at a rate of ω0 /Q. This suggests that the
energy decays most rapidly in the vicinity of critical damping or Q ∼ 21 .
Exercise 2.2.1 A 0.5 kg mass hangs on the end of a 1.0-m-long pendulum that is part
of a grandfather clock. In ten minutes, the amplitude of the pendulum decays from
4.0 cm to 0.5 mm. (a) What is the Q of this pendulum? (b) Show that the fraction of
energy lost each cycle is approximately 2π/Q.
Exercise 2.2.2 A church bell produces a sound at 220 Hz (the sound of the musical
note A below middle C). After it is struck, the sound dies down in about 5 s. What is
the approximate Q of the bell?
Exercise 2.2.3 For a very weakly damped oscillating pendulum with amplitude z0 ,
show (i) that the total energy is given by ET ' 21 mω20 z20 , (ii) that the energy dissipated
in one cycle is W f ' 12 Γω20 z20 (2π/ω0 ), and (iii) that the ratio of energy stored to en-
ergy lost in one cycle is ET /W f ' ω0 /2πγ = Q/2π. This last equation Q = 2πET /W f
is often used as the definition of Q. It’s worth remembering.
2.3 Shock absorbers
Perhaps the most common use of damped oscillators is to mitigate the potentially delete-
rious effects of a sudden shock to a system. A sensitive piece of electronics like a mobile
phone is more likely to break if it falls one the floor unprotected than if has an elastic rub-
ber coating to cushion the impact. We don’t want the cushion to be too springy, however,
as it is generally preferable for the phone to bounce only once or twice before coming to
rest. The “springiness” of the protective rubber coating should be appropriately damped.
Another example of how damped oscillators are used is found in the shock absorbers of an
automobile. The idea is for the shock absorbers to cushion the impact to the car (and its
passengers) when it runs over a bump or hole in the road by mounting springs between the
wheels and the car frame. However, we would not want the springs to cause running over
a bump to send the car into large persistent oscillations, which could cause loss of control
(of either the car or the passengers’ stomachs). Any oscillations that are caused by the sus-
pension springs should die out quickly after the springs have done their job of cushioning
the impact.
t
Let’s consider a simple model of a car suspension system. The suspension system con-
sists of a heavy duty spring and a viscous damper, which we model as a “dashpot," as
shown in Fig. 2.5. The spring connects the car frame (green) to axle (black dot) of the
wheel (large gray circle). The dashpot, which acts as a damping mechanism, consists of
(1) hollow cylinder with a closed bottom (red) that is connected to the axle and (2) a piston
(blue) that fits inside the cylinder and is connected directly to the car frame. When the
car frame moves relative to the axle, the spring either stretches or compresses while the
dashpot provides the viscous friction force proportional to the velocity of piston (and the
car frame) relative to the cylinder (and the wheel axle). The distance from the axle to the
car frame is z. We take z = z0 to be the resting distance between the car frame and the axle.
In this state the net force F on the car is zero: F = −k(z0 − zu ) − Mg = 0, where k is the
spring constant, zu corresponds to the position of the unstretched spring, and M is the mass
of the car that is supported by this spring. When the car moves vertically relative to the
axle, the viscous damping force provided by the dashpot is given by Fd = −Γż, where ż is
the vertical velocity of the car and Γ is the viscous damping coefficient, which is taken to
be constant.
Putting this all together, the equation of motion for the vertical movement of the car is
F = −k(z − zu ) − Mg − Γż = Mz̈ . (2.69)
Defining ξ = z − z0 as the distance from the equilibrium position z0 and using the equilib-
rium condition that k(zu − z0 ) = Mg, Eq. (2.69) becomes
−kξ − Γ ξ̇ = M ξ̈ . (2.70)
Rearranging terms, this becomes
ξ̈ + γξ̇ + ω20 ξ = 0 (2.71)
√
where γ = Γ/M and ω20 = k/M. This is the same equation of motion we previously
encountered for the damped pendulum (see Eq. (2.4)). Notice the trick we used above to
t
Fig. 2.5 Suspension system with shock absorber on an automobile wheel. Here, the dashpot
(red and blue) fits inside the hollow of the spring (black).
45 Shock absorbers
t
get rid of the constant (inhomogeneous) term involving zu in Eq. (2.69). By defining the
zero of the coordinate ξ to be at the equilibrium position of the spring under the weight of
the car, and using the equilibrium condition that k(zu − z0 ) = Mg, we obtained Eq. (2.71),
an equation of motion without a constant term. This is the same trick we introduced in
§1.3.2. It’s a good trick to remember!
Let’s use our results to design the suspension for a car. We suppose that one wheel
supports about 400 kg or about 1/4 of the mass of a 1600 kg car. Let’s suppose that the car
is moving at about 3 x = 15 m/s (54 kph ' 34 mph) when it runs over a sudden drop in the
road of d = 5 cm. For the sake of simplicity we shall assume that the distance between the
car frame and the axle remains fixed at z0 during the fall. We further assume that when the
tire hits the ground it abruptly stops its vertical motion while the car frame initially keeps
moving. Its initial vertical velocity is 3z = (2gd)1/2 ' 1.0 m/s downward. Thus, the initial
conditions are ξ̇ = −1.0 m/s and ξ = 0 (i.e. z = z0 ).
To get a sense of what is going on, we start by assuming that we want the suspension
to be critically damped. The boundary conditions at t = 0 are zero displacement and a
nonzero velocity ẋii . Our equations of motion are identical to those for the damped spring,
Eq. (2.44), with ξ(t) replacing z(t):
ξc (t) = ξ̇i t e−γt/2 . (2.72)
In designing the shock absorber, there are at least two important considerations. First, we
do not want the shock of the car hitting the ground after the drop to be too large. Second,
we want the car to recover quickly from the effects of the drop in the road. Let’s consider
the shock first. Our concern is simply that the vertical acceleration associated with the car
hitting the ground not be too big. Differentiating Eq. (2.72) twice to obtain the vertical
acceleration of the car reveals that the maximum acceleration amax occurs at t = 0 and
is given by amax = γξ̇i , that is, the damping constant times the initial vertical velocity
when the wheel his the ground. Since this is a modest size drop, let’s see if we can limit
amax to 12 g ' 5 m/s2 . (For reference, the accelerations experienced while walking have an
amplitude of about 15 g.) With ξ̇ = −1.0 m/s, this gives γ = 5 s−1 .
We also want the car response and recovery time to occur on a time scale no longer
than a second or so. There are two time scales in the problem, the damping time, which
is given by 2/γ, and the inverse frequency (or period) given by ω−1 0 = 1/Qγ. We expect
the optimal solution to be near critical damping, for which Q = 21 . Thus, both time scales
2/γ and ω−1 0 = 1/Qγ are on the order of 0.4 s. Since the maximum acceleration is set by
the damping constant, we will hold γ fixed and vary ω0 to vary Q near its value for the
critically damped case.
In Fig. 2.6 we plot the response given by Eq. (2.72). For the critically damped case,
the spring compresses (ξ < 0), reaches a maximum compression of about 15 cm at about
t = 0.4 s, and then relaxes back to zero displacement after about 3 s.
To get a better sense of how to optimize our design, we examine what happens for the
slightly underdamped and overdamped cases. Once again we can adapt our results from
the damped oscillator to get the responses for the underdamped and overdamped cases.
From Eqs. (2.26) and (2.37), we obtain the following results for the underdamped and
t
overdamped cases respectively
ξ̇i −γt/2
ξu (t) = e sin ωu t (2.73)
ωu
ξ̇i
ξo (t) = e−γt/2 sinh βt (2.74)
β
where ωu is given by Eq. (2.16) and β is given by Eq. (2.31). In Fig. 2.6. Returning to our
design criterion that the maximum acceleration be given by amax ≈ 5 m/s2 , we once again
calculate amax from Eqs. (2.26) and (2.37) and find that amax = γξ̇i for both the under-
damped and overdamped cases, just as we found for the critically damped case. Therefore,
we shall keep the same value of γ = 5 s−1 . For the underdamped and overdamped cases,
there is another time scale that we can tune, namely ω−10 . For the critically damped case, ω0
is fixed since Q = ω0 /γ = 12 , but for the other cases it can be varied since Q = ω0 /γ > 12
for the underdamped case and Q = ω0 /γ < 12 for the overdamped case. Because we expect
that optimal damping occurs near critical damping, we vary Q by only a factor of two in
ξ (t) = z(t) − z0 (m)
0.0 t (s)
1 2 3 4
Q=1
Q = 1/2
−0.1
Q = 1/4
−0.2
a(t) (m/s2 )
5
0 t (s)
1 2 3 4
t
Fig. 2.6 Response of automobile shock absorbers to running over a pothole. Top:
displacement from equilibrium wheel position. Bottom: vertical acceleration of
automobile. The damping is the same for each curve; Q is varied by changing the
spring constant k = Mγ2 Q2 .
47 Shock absorbers
t
either direction choosing Q = 1 and Q = 14 for the underdamped and overdamped cases,
respectively.
The acceleration profiles for the three cases are very similar, with the critically damped
and overdamped cases being nearly indistinguishable (see bottom of Fig. 2.6). On the other
hand, the displacement curves for the three cases differ significantly from each other. For
the overdamped case, the wheel is still more than 5 cm from its equilibrium position 4 s
after impact. While such a long relaxation time might have been expected, it is interesting
that the wheel returns to its equilibrium position considerably more rapidly for the under-
damped than for the critically damped case. By increasing Q while keeping γ constant for
the underdamped case, we have increased the spring constant k, since k = Mω20 = Mγ2 Q2 .
The increased value of k limits the maximum displacement, which explains why the maxi-
mum displacement is smallest for the underdamped case. Of course, the underdamping also
causes the overshoot as the wheel approaches it equilibrium position. Increasing Q > 1
would increase the number of oscillations and thus defeat the damping function of the
shock absorber. Nevertheless, the curves in Fig. 2.6 suggest that a small degree of un-
derdamping, which implies a stiffer spring, is desirable in the design of shock absorbers.
However, it is clear that the price paid for returning to equilibrium quicker is a somewhat
stiffer ride, as indicated by somewhat greater average absolute acceleration evident in Fig.
2.6. This trade off between a softer ride but slower response and a stiffer ride but quicker
response is one reason why the family sedan has a more comfortable ride than a sports car,
but the sports car gives the driver more control.
For our final design, we therefore set γ = 5 s−1 and Q = 1. For M = 400 kg, this
gives a spring constant of k = Mω20 = Mγ2 Q2 = 10, 000 N/m and a damping constant of
Γ = Mγ = 2, 000 N · s/m. Real cars have significantly stiffer springs. One reason is that
the air-filled tires take up much of the shock of driving over a bump, so the shock absorbers
only have to do only a part of the job.
t
Problems
2.1 Consider a pendulum consisting of a mass m = 0.20 kg at the end of a massless rod
1.20 m long attached to the ceiling. The mass is displaced 5.0 cm from equilibrium,
let go, and allowed to oscillate freely. After 10 oscillations, its amplitude is 4.0 cm.
(a) Is the system overdamped, underdamped, or critically damped? Is the value of Q
much greater than 1 or much less than 1 (do not calculate Q yet). Briefly explain
(1 or 2 short sentences!).
(b) Using the fact that the amplitude of oscillations decays approximately as e−γt/2 ,
determine the approximate numerical value (correct to within 5%) of Q for this
system.
2.2 In this problem you will show that the general solution for the critically damped
oscillator is given by Eq. (2.43), which for your convenience we rewrite here
zc (t) = (c1 + c2 t) e−γt/2 . (2.75)
Start with the solution for the overdamped case, Eq. (2.32), written below in a slightly
more convenient form
z(t) = e−γt/2 b1 e−βt + b2 eβt ,

h i
(2.76)
q
where β = (γ/2)2 − ω20 . The strategy is to find z(t) from Eq. (2.76) in the limit
β → 0, which corresponds to the critically damped case. Since you can make β
arbitrarily small, you can expand the exponentials e−βt and eβt in a power series about
βt = 0. Go ahead and do this. Then gather together all the even and odd terms of the
power two power series and show that can you obtain
(βt)2 (βt)4
" !
z(t) = e−γt/2
(b1 + b2 ) 1 + + + ...
2! 4!
(βt)2 (βt)4
!#
+ (b2 − b1 )βt 1 + + + ... . (2.77)
3! 5!
Letting c1 = (b1 + b2 ) and c2 = (b2 − b1 )β, you have the solution given by Eq. (2.43).
How is it that the terms in the power series go to zero as β → 0 but the term (b2 −b1 )β
can survive and remain finite?
2.3 Find the solutions for the damped oscillator for the initial conditions that z(t = 0) =
z0 > 0 and ż(t = 0) = 0. In particular,
(a) starting from Eq. (2.20), show that
γ
!
zu (t) = z0 e
−γt/2
cos ωu t + sin ωu t (2.78)
2ωu
for the underdamped case.
49 Problems
t
(b) Starting from Eq. (2.32), show that
γ
!
zo (t) = z0 e
−γt/2
cosh βt + sinh βt (2.79)
2β
for the overdamped case.
(c) Starting from Eq. (2.43), show that
γt −γt/2
zc (t) = z0 1 + e (2.80)
2
for the critically damped case.
(d) On a single graph, plot of z(t) vs. t for Q = { 14 , 12 , 1, 2, 4} and z0 = 0.05 over
the range 0 ≤ t ≤ 5. For the purpose of your plots set the period of oscillation
T = 1 ⇒ ω0 = 2π.
(e) Which of your plots are overdamped, underdamped, or critically damped. For
any underdamped cases, how much is the period changed from what it would
have been if there were no damping?
2.4 According to classical electrodynamics, an accelerating charge radiates energy at a
rate of Kq2 a2 /c3 where a is its acceleration, q is the charge, c is the speed of light,
and K = 6.0 × 109 N m2 /C2 . Suppose that the electron is oscillating during during
one cycle of its motion according to z = A sin ωt.
(a) Show that the energy radiated during one cycle is given by πKe2 ω3 A2 /c3 , where
e is the electronic charge.
(b) Recalling that the energy of a simple harmonic oscillator is 21 mω2 A2 , show that
Q = mc3 /Ke2 ω.
(c) Using a typical value of ω for visible light, estimate the “lifetime" of the radiating
system.
2.5 Redo the problem in §2.3 of a car going over a drop in the road but relax the assump-
tion that the spring does not extend while the car is falling. Use the same parameters
used in §2.3. Take the radius and mass of the car wheel to be 0.225 m and 8 kg, re-
spectively. Determine the time-dependent displacement ξ and acceleration. Present
your results in graphical format similar to that of Fig. 2.6.
2.6 A circular metal plate with moment of inertia I = 0.02 kg-m2 is suspended from a
metal wire. When the plate is twisted through an angle θ away from its equilibrium
position, it experiences a torque −kθ. If the plate rotates between two magnet heads,
it experiences a dissipative torque −bθ̇. Thus the equation of motion for the plate in
the presence of both the magnets and an external torque τext (t) is
I θ̈ + kθ + bθ̇ = τext (t) .
When the magnet heads are removed, there is no dissipative torque, and the system
undergoes simple harmonic motion with period T = 0.5 s. When the magnet heads
are in place, the system has a quality factor Q = 8.
(a) What is the torque constant k? Be sure to specify units. (Ans: k = 0.32π2 kg-m2 /s2 ))
(b) What is the damping constant γ? Be sure to specifyunits. (Ans: γ = 0.01π kg-m2 /s)
t
t
Fig. 2.7 Torsion Oscillator
(c) At time t = 0, the plate is given an impulse such that its initial angular velocity
is π radians/sec. What is the angular displacement θ after 0.125 seconds?
2.7 A mass-and-spring harmonic oscillator oscillates at a frequency ω0 = 2π/T 0 when it
is undamped. It is placed in a viscous medium, and it acquires a quality factor Q = 1.
(a) At time t = 0 the mass, at rest at its equilibrium position at the origin, is de-
livered at impulse that gives it a velocity 30 . Determine the time T 1 at which
the mass first returns to the origin. What is the position z2 of the mass √ at time
T 1 /2?. Express
√ your answer
√ in terms of T 0 and 30 . (Ans: T 1 = T 0 / 3, z2 =
[30 T 0 /( 3π)] exp[−π/(2 3])
(b) The oscillator is immersed in another viscous medium, and its quality factor
decreases to 1/4. If the particle is displaced a distance z0 from the origin at time
t = 0 and released, what is its position at time T 0 ? Ans:
z0
z(T 0 ) = √ α+ e−2πα− − α− e−2πα+ ,
2 3
√
where α± = 2 ± 3
2.8 (French 3-14) An object of mass 0.2 kg is hung from a spring whose spring constant
k is 80 N/m. The object is subject to a resistive force given by −Γ3, where 3 is the
velocity in m/s.
(a) Set up the differential equation
√ of motion for the free oscillation of the system.
(b) If the damped frequency is 3/2 of the undamped frequency, what is the value
of the constant Γ? [Ans: Γ = 4 kg/s]
(c) What is the Q of the system, and by what factor is the amplitude of√oscillation
reduced after 10 complete cycles? [Ans: Q = 1, A(10)/A(0) = e−20π/ 3 = 1.76 ×
10−16 .]
2.9 (French 3-15) Many oscillatory systems, although the loss or dissipation mechanism
is not analogous to viscous damping, show an exponential decrease in their stored
average energy with time E = E 0 e−γt . A Q for such oscillators may be defined by
using the definition Q = ω0 /γ, where ω0 is the natural angular frequency.
(a) When the note “middle C" on the piano is struck, its energy of oscillation de-
creases to one half its initial value in about 1 sec. The frequency of middle C is
256 Hz. What is the Q of the system? [Ans: Q = 2πν/ ln 2 = 2321]
51 Problems
t
(b) If the note an octave higher (512 Hz) takes about the same time for its energy to
decay, what is its Q. [Ans: Q = 4642]
(c) A free, damped harmonic oscillator, consisting of a mass of m = 0.1 kg moving
in a viscous fluid of damping coefficient Γ (Fviscous = −Γ3), and attached to
a spring of spring constant k = 0.9 N/m, is observed as it performs damped
oscillatory motion. Its average energy decays to 1/e of its initial value in about
4 sec. What is the Q of the oscillator? What is the value of Γ? [Ans: Q = 12,
Γ = 0.025 kg/s]
2.10 Consider a damped harmonic oscillator with a natural frequency ω0 and inverse
damping time γ whose displacement as a function of time is given by z(t).
(a) Derive an expression for z(t) for the over damped case subject to initial condi-
tions z(0) = z0 and ż(0) ≡ 30 . Ans:
1 sinh βt
" #
z(t) = e −γt/2
z0 cosh βt + (30 + γz0 ) .
2 β
q
where β = (γ/2)2 − ω20 .
(b) Derive the corresponding expressions for z(t) for the critically damped and under-
damped oscillators using only the expression for the over-damped oscillator
q the facts that β = 0 in the critically damped oscillator and β → iωu =
and
i ω20 − (γ/2)2 in the under-damped oscillator. Verify that z(0) = z0 and 3(0) = 30
in both cases. Ans:
1
" #
z(t) = e −γt/2
z0 + (30 + γz0 )t critical,
2
1 sin ωu t
" #
z(t) = e −γt/2
z0 cos ωu t + (30 + γz0 ) . under-damped
2 ωu
3 Resonance
The pendulum in a well-designed mechanical clock is an exquisite time-keeping device. A

good pendulum clock typically has an accuracy of about 1 part in 105 , or about 10 seconds
a month. However, a lone pendulum without any other mechanism is a poor time-keeping
device because the oscillations that measure the passage of time gradually die out as fric-
tion saps energy from each successive swing. The only way to keep a such a clock going
is to periodically inject some energy. This is how the pendulum on a mechanical clock
sustains its motion. By means of a mechanism involving a ratchet and a falling weight or
a compressed spring, the pendulum is given a small impulse each cycle. A key feature that
permits the impulse to be so small is that it is synchronous with the natural frequency of the
pendulum. That is, the impulse has the same frequency ω0 as the natural frequency of the
pendulum. This phenomenon, that periodic forcing near the natural frequency of an oscil-
lator results in oscillations with much larger amplitudes than forcing at other frequencies,
is called resonance.
Resonance is widely observed in nature and exploited by technology. One of the princi-
pal uses of resonance is the detection of waves, most notably sound and electromagnetic
waves, including radio, light, and z-rays. For example, radio waves, a vast spectrum of
which fills almost every space we occupy every hour of every day, can be detected when
they set into motion electrons in an antenna. The radio waves create currents that oscillate
over an incredibly broad range of frequencies, from several hundred kilohertz to hundreds
of megahertz. How is it that the radio tuner finds the signal of that one radio station you
want to listen to? The answer lies in the phenomenon of resonance.
The aim of this chapter is to provide you with a quantitative introduction to oscillations
and resonance. By the end of the chapter, you should understand how a radio receiver
tunes in a specific frequency broadcast, how a seismometer detects a broad range of seis-
mic movements, and how a pendulum works. And so we proceed, from the simple to the
sublime.
3.1 Forced oscillators and resonance
As noted above, the only way to maintain the motion of an oscillator indefinitely is to
periodically inject some energy. One example is a pendulum clock, which was previously
discussed. Another example is a child on a swing with a parent periodically giving the
child a push. The pushing is most effective if it occurs in sync with the natural period of
52
53 Forced oscillators and resonance
t
the swing and if it is in the proper direction. That is, the parent should push the swing in
the direction it is already going, and certainly not in the opposite direction, which would
slow down the swing. For the forcing to be most effective, therefore, it should occur in
sync with the natural frequency of the oscillator and it must be properly phased. We will
look for these features when we obtain solutions to the equation of motion for a forced
oscillator.
3.1.1 Resonance
We begin our study of resonance by introducing periodic forcing to a damped oscillator.

Here we consider a weight of mass m suspended by a spring of spring constant k and sub-
jected to a friction force proportional the weight’s velocity 3, Fd = −Γ3, which can arise,
for example, from a surrounding viscous fluid or from confinement to a hollow cylinder
providing lubricated contact between the weight and the cylinder as depicted in Fig. 3.1.
The piston is driven by a sinusoidally oscillating force F(t) = F0 cos ωt, which might be
provided by coupling magnetically to a metallic weight.
Measurement of the time-dependent position z(t) of the weight as a function of ω and
F0 reveals the following properties:
(a) After an initial period after the driving force is turned on, the system reaches a steady
state in which the vertical position of the weight undergoes simple harmonic motion
at the same frequency ω as the driving force but with a phase lag, φ(ω), relative to the
driving force: y(t) = y0 + A(ω) cos(ωt −φ), where y0 = mg/k is the equilibrium position
of the weight and A(ω) is the maximum amplitude.
√
(b) When ω ω0 = k/m, the driving force and the weight displacement follow each
other (φ(ω) ≈ 0): they reach maxima and minima at nearly the same time. The ampli-
tude A(ω) increases with increasing ω.
(c) When ω ω0 , the driving force and displacement are almost completely out of phase,
φ ≈ π: the driving force is at a maximum when the displacement is at a maximum and
vice versa. The displacement amplitude decreases with increasing frequency.
t
Fig. 3.1 Mass in a hollow cylinder suspended by a spring and driven by an oscillating
magnetic field from coil. The motion is damped by the lubricated contact between the
mass and the cylinder.
54 Resonance
t
(d) The amplitude A(ω) reaches a maximum in the vicinity of ω = ω0 with a peak value
that is higher the smaller the damping (reduced for example by switching from a vis-
cous fluid like glycerol to a less viscous one like water). The peak also gets narrower
as damping is decreased. This is resonance. The phase displacement-force phase lag is
π/2 at ω = ω0 : displacement is zero when the driving force is a maximum.
(e) The power required to drive the system follows a curve as a function of ω that is similar
to that of the amplitude with a peak at ω0 and dying off with ω at large ω, but going to
zero at ω = 0.
We now show how all of these properties follow from the harmonic-oscillator equation
of motion in the presence of a sinusoidal driving force. Gravity acts downward on the
weight in the geometry of Fig. 3.1. Following the procedure discussed in §1.3.2, we define
z(t) = y(t) − mg/k to remove gravity from the equation of motion, to obtain
d2 z dz
m = −kz − Γ + F0 cos ωt . (3.1)
dt2 dt
After some rearrangement of terms, this equation can be rewritten as
d2 z dz F0
2
+ γ + ω20 z = cos ωt , (3.2)
dt dt m
√
where ω0 ≡ k/m and γ ≡ Γ/m. Equation (3.2) has the same form as Eq. (2.4), the
equation for the unforced damped oscillator, except that Eq. (3.2) includes a forcing term,
which is an explicit function of the independent variable t. Thus, Eq. (3.2) is inhomoge-
neous, whereas Eq. (2.4) is homogeneous, as discussed in §1.1.
We look for solutions to Eq. (3.2) using the same approach we used for the unforced
pendulum, by assuming a solution of the form z(t) = Ãe−iωt , where Ã is complex. To take
advantage of the algebraic simplicity that exponential notation affords, we write the forcing
term as
F0 cos ωt = F0 Re e−iωt , (3.3)
taking only the real part of e−iωt , which is cos ωt. Thus, Eq. (3.2) becomes
z̈ + γż + ω20 z = (F0 /m) e−iωt , (3.4)
where we have omitted explicitly writing Re for notational simplicity. We just have to
remember to take only the real part of any solution we obtain at the end of our calculation.
Initially, we looking only for steady state solutions, that is, for solutions that have only
a simple sinusoidal time dependence of the form
z(t) = Ãe−iω t ,
0
(3.5)
where we take the weight to oscillate at a frequency ω0 to allow for the possibility that the
weight oscillates at a frequency different from either the natural frequency ω0 or the drive
frequency ω. Keep in mind that Ã and e−iω t are complex, which allows for there to be a
0
phase difference between the driving force and the response of the oscillator. In the end we
will keep only the real part of z(t). The other terms in the equation, including ω0 , γ, m, and
F0 are real constants that describe the oscillator and the nature of the forcing.
t
To look for solutions, we substitute Eq. (3.5) into Eq. (3.4), which gives
−ω0 2 − iγω0 + ω20 Ã e−iω t = (F0 /m) e−iωt .

0

(3.6)
Multiplying both sides of the equation by eiω t yields

0
−ω0 2 − iγω0 + ω20 Ã = (F0 /m) e−i(ω−ω )t .

0

(3.7)
Notice that all the time dependence in this equation is on the right hand side. The left
hand side is completely independent of time! For the two sides of the equation to be equal,
the right hand side must be independent of time as well. The only way this can happen is
if ω0 = ω. Thus, even before finding a full solution to the problem, we arrive at a very
important result: If a linear oscillator is driven at a frequency ω, then the oscillator will
respond only at that frequency. This means that it won’t respond at the natural frequency
ω0 . Now this turns out to be true only in steady state, that is, for times much longer than
the natural damping time of the oscillator, which from the previous section is typically of
the order of γ−1 = Q/ω0 , where Q is the quality factor introduced in §2.1.1. It turns out
that we have left out the transient solution in our work above. Don’t worry, we will put it
in shortly. Nevertheless, the point is important enough to reiterate. If a linear oscillator is
driven at a frequency ω, its steady state response will be only at that frequency.
With ω0 = ω, our trial solution becomes
z(t) = Ã(ω)e−iωt , (3.8)
where we display explicitly that Ã depends on frequency. We will not always do so in what
follows. Proceeding with this trial solution, Eq. (3.7) becomes
−ω2 − iγ ω + ω20 Ã = F0 /m .

(3.9)
Solving for Ã gives

F0 /m
Ã(ω) = . (3.10)
ω20 − ω2 − iγ ω
Writing the amplitude in polar form
Ã(ω) = A0 (ω)eiφ(ω) , (3.11)
where A0 (ω) and φ(ω) are the real amplitude and phase, Eq. (3.10) yields 1
F0 /m
A0 (ω) = q 2 (3.12)
ω20 − ω2 + γ2 ω2
 γω 
 
φ(ω) = tan−1  2  . (3.13)
ω0 − ω2
Writing the complex amplitude as A = A0 eiφ , Eq. (3.8) becomes
z(t) = Ãe−iωt = A0 e−i(ωt−φ) . (3.14)

√
1 It’s worth pausing here to prove to yourself that 1/(a − ib) = eiφ / a2 + b2 , where tan φ = b/a.
56 Resonance
t
A 0 (ω )
15
Q = 16
10
√
∆ω ≈ 3 γ ∼ Q−1
0 ω /ω0
0 1 2 3
φ (ω )
π
Q = 16
π /2 Q = 1/2
0 ω /ω0
0 1 2 3
t
Fig. 3.2 Amplitude (top) and phase of oscillator response z(t) to periodic forcing of f0 cos ωt for
values of the quality factor Q from the lightest to the darkest curves: 12 , 1, 2, 4, 8, 16.
The height of the peak of the resonance curve (top) is given by a0 (ω0 ) = Q( f0 /ω20 ) and
width ∆ω is proportional to 1/Q in the limit that Q 1. Phase (bottom) of the oscillator
response for values of Q ranging from 1 to 16, with shadings corresponding to the top
curves. The curves corresponding to Q = 1 are shown as dashed lines.
Retaining only the real part, we obtain
z(t) = Re Ã(ω) e−iωt = A0 cos(ωt − φ) (3.15a)

F0 /m
= q 2 cos(ωt − φ) , (3.15b)
ω20 − ω2 + γ2 ω2
where the phase φ is given by Eq. (3.13). Equation (3.15) represents the complete solution
to the steady state driven damped oscillator. This important result states that amplitude and
the phase of the oscillator are functions of the the drive frequency ω, the natural frequency
ω0 , and the damping rate γ.
Equation (3.15) provides the steady state driven solution for z(t) when the physical driv-
ing force is F0 cos ωt. What happens when the physical driving force is F0 sin ωt or even
F0 cos(ωt − β) for arbitrary angle β? One only has to replace ωt by ωt − β, with β = π/2
for the F0 sin ωt case, in z(t). Thus if for β = π/2, z(t) = Ã sin(ωt − φ).
t
Figure 3.2 shows the amplitude and phase of the oscillator given by Eqs. (3.12) and
(3.13). The most striking feature of the plots is the dramatic increase of the amplitude A0
of the oscillator when the driving frequency ω is near the natural oscillation frequency ω0
for the case of weak damping, that is, when Q 1. This dramatic increase in the response
of the oscillator for driving frequencies near ω0 is known as resonance and is a universal
feature of high-Q oscillators.
The frequency at which the amplitude is a maximum is given by
s
1
ωmax = ω0 1 − , (3.16)
2Q2
which, for Q 1, is very nearly equation to the free (or natural) frequency of oscillation
ω0 , which is often referred to as the resonant frequency.
For Q 1, the maximum amplitude of the oscillator is given by
F0 F0
A0 (ωmax ) ' A0 (ω0 ) = Q =Q , (3.17)
mω20 k
which means that for Q 1, the peak amplitude is enhanced by a factor of Q over the the
displacement F0 /k that would be expected for a constant (non-oscillatory) force F0 .
The range of frequencies around the resonant frequency where the amplitude is large
is determined by the damping
√ rate γ. The amplitude is reduced to half its peak value at
frequencies ω ' ω0 ± ( 3/2)γ. The so-called full width at half maximum (fwhm), which
is defined as the width ∆ω of the resonant curve at half of its maximum amplitude A(ω0 ) is
given by
√ ω0
∆ω ≈ 3 γ ∼ , (3.18)
Q
The quality factor Q thus provides a convenient quantitative measure of the resonant curve:
the peak height increases in proportional to Q and the width narrows in proportion to Q−1 .
Exercise 3.1.1 Starting from Eq. (3.12), show that the for Q 1, the full width half
maximum frequency is given by Eq. (3.18). Hint: In the large Q limit where ω0 γ,
you can approximate the terms in the denominator of Eq. (3.12) using
ω20 − ω2 = (ω0 − ω)(ω0 + ω) ≈ 12 ∆ω (2ω0 ) = ∆ω ω0

(3.19)
γω ≈ γω0 . (3.20)
Justify these and any other approximations you make in arriving at Eq. (3.18).
The phase φ(ω) of the oscillator, defined in Eq. (3.15a) and given by Eq. (3.13), also
varies with the frequency ω, especially in the vicinity of ω = ω0 . For frequencies much
smaller than the resonant frequency, φ is nearly zero, meaning that the response of the
oscillator is in phase with the forcing. This makes sense: if you push very slowly on a os-
cillator, the applied force is simply balanced by the spring restoring force. Indeed, because
the acceleration and velocity are very small in this case, the first two terms in Eq. (3.4) are
very small compared to the third (see Eq. (3.9)). Thus, the equation of motion becomes
z(t) ' (F0 /mω20 ) cos ωt = (F0 /k) cos ωt, which is consistent with Eq. (3.15b) in the limit
58 Resonance
t
z(t) F(t)
5 ω = 21 ω0
0 t/T0
1 2 3
−5
z(t) F(t)
5 ω = ω0
0 t/T0
1 2 3
−5
z(t) F(t)
5 ω = 2ω0
0 t/T0
1 2 3
−5
t
Fig. 3.3 Motion z(t) of oscillator (black lines) in response to periodic forcing F(t) (dashed gray
lines) for (a) ω = ω0 /2 (below resonance), (b) ω = ω0 (at resonance), and (c) ω = 2ω0
(above resonance). In all cases, the forcing is given by F(t) = F0 cos ωt. Time is
measured in units of the natural period T 0 = 2π/ω0 and Q = 8 for all curves.
that ω ω0 . In this case the amplitude is just the equilibrium displacement F0 /k asso-
ciated with a static force of F0 . The driving force and oscillator response in this limit are
illustrated in Fig. 3.3(a).
As the frequency ω approaches the resonant frequency ω0 , φ(ω) increases and the re-
sponse of the oscillator starts to lag behind the driving force. That is, the oscillator reaches
its peak a short time after the forcing goes through its peak value. Precisely at resonance
when ω = ω0 , φ(ω0 ) = π/2 and the oscillator position lags the driving phase of the driv-
ing force by 90◦ . On the other hand, the velocity ż is in phase with the driving force at
resonance. Thus, the force is always pushing in a direction that tends to enhance the oscil-
lator’s motion. This is consistent with your experience of pushing someone on a swing; you
push them in the same direction they are already moving. The driving force and oscillator
response at resonance are illustrated in Fig. 3.3(b).
As the driving frequency increases, the oscillator can no longer keep up with the driving
force. When this occurs, the driving force actually opposes the motion of the oscillator for
part of its cycle, resulting in a loss in amplitude. For ω ω0 , the driving force is 90◦ out
of phase with the velocity so that the force actually opposes and enhances the oscillator
motion for nearly equal amounts of time resulting in a vanishingly small amplitude. In this
t
15 Q = 16
10
5 χ′
χ ′′
0 ω /ω0
1 2 3
15
−10
t
Fig. 3.4 Response functions of an oscillator: in phase χ0 (solid line) and out-of-phase χ00
(dashed line).
case, the displacement and driving force are almost completely out of phase. This situation
is illustrated in Fig. 3.3(c).
We end this section by summarizing the response of the driven oscillator, both in ampli-
tude and phase, at the resonant frequency ω0 and for frequencies much less than and much
greater than ω0 :
(F0 /mω20 ) = F0 /k for ω → 0 0 for ω → 0
 


 


 
2
A0 (ω) '  Q (F0 /mω0 ) = QF0 /k for ω → ω0 , φ(ω) '  π/2 for ω → ω0 . (3.21)

 


 

2
(F0 /mω ) → 0 for ω → ∞ π for ω → ∞

 

3.1.2 Response functions
In the previous section we characterized the response of a simple harmonic oscillator to

an external periodic driving force in terms of the amplitude and phase of the oscillator.
An alternate way to express this same information about the response of the system is in
terms of the in-phase and out-of-phase amplitudes of the motion. One reason why this
is interesting is that there is a close relationship between the energy dissipated by the
oscillator and its out-of-phase amplitude.
For an oscillating driving force of F(t) = F0 cos ωt, we found we found the solution
z(t) = Ãe−iωt , where Ã and e−iωt are both complex. Writing the amplitude in terms of its
real and imaginary parts Ã = AR + iAI , the displacement of the oscillator becomes
n o
z(t) = Re A e−iωt (3.22)
= Re {(AR + iAI ) (cos ωt − i sin ωt)} (3.23)
= AR cos ωt + AI sin ωt , (3.24)
where we have kept only the real part of the expression. The cos ωt term is, of course, in
60 Resonance
t
phase and the sin ωt is out of phase with the driving force F0 cos ωt. Thus, from Eq. (3.24)
we see that amplitudes of the in-phase and out-of-phase response are given by AR and
AI the real and imaginary parts of A. With this in mind we define a frequency-dependent
response function for each part
AR (ω) 1 ω20 − ω2
χ0 (ω) = = (3.25)
F0 m (ω20 − ω2 )2 + (γ ω)2
AI (ω) 1 γω
χ00 (ω) = = . (3.26)
F0 m (ω20 − ω2 )2 + (γ ω)2
and the complex response function

A(ω)
χ(ω) = = χ0 (ω) + iχ00 (ω) , (3.27)
F0
where χ0 (ω) and χ00 (ω) are its real and imaginary parts. It is clear that
χ(ω) = |χ(ω)|eiφ(ω) , (3.28)
where |χ(ω)| = A0 (ω)/F0 = (χ00 )2 + (χ00 )2 and φ(ω) = tan−1 (χ00 /χ0 ) is equal to phase of
p
Ã of Eq. (3.13).
Figure 3.4 shows the frequency-dependent in-phase and out-of-phase response functions
χ0 (ω) and χ00 (ω), respectively. The limiting forms of these functions are are summarized
as
1/k for ω → 0 0 for ω → 0
 


 


 
χ (ω) ' 
0
0 for ω → ω0 , χ (ω) ' 
00
Q/k for ω → ω0 . (3.29)

 


 

−1/(mω2 ) → 0 for ω → ∞

 0


for ω → ∞
A number of properties of χ(ω) should be noted. χ0 (ω) is even under ω → −ω, positive for
ω < ω0 , equal to zero at ω0 = 0, and negative for ω > ω0 . χ00 (ω), on the other hand, is odd
under ω → −ω and thus equal to zero at ω = 0, positive for all ω > 0, peaked, like A0 (ω)
and |χ(ω)|, in the vicinity of ω = ω0 with a width that goes to zero with γ, and equal to Q/k
at ω = ω0 . The limiting forms of Eq. (3.29) can be deduced directly from the equations
of motion, Eqs. (3.2) and (3.9). At ω = 0, the inertial (mz̈) and dissipative (γż) terms are
zero, and all that remains on the right-hand side of the equation is the static term kz, and
χ(0) = 1/k; when ω → ∞, the inertial term dominates so that χ(ω → ∞) = −1/(mω2 );
finally at ω = ω0 , only the dissipative term remains, and χ(ω = ω0 ) = i/(mω0 γ), implying
that the response is purely imaginary in this limit and that, as a result, φ = π/2.
In the current context, the equations of motion are linear in both z(t) and F(t), and χ(ω)
is simply Ã(ω)/F0 . In most problems, however, Ã or its analog depend nonlinearly on F(t),
but the response function χ(ω) is still well defined: it is just the derivative of Ã with respect
to F0 evaluated at F0 = 0. Thus, χ(ω) is just a generalization to non-zero frequency of a
static susceptibility, like the magnetic susceptibility, ∂M/∂H, where M is the magnetization
and H is the magnetic field, or the electric susceptibility ∂P/∂E, where P is the electric
polarization and E is the electric field. Thus, we expect χ(ω = 0) to be the “zero-frequency
susceptibility” dx/dF0 = 1/k as it is. As we shall see in a later chapter, χ(ω) is the Fourier
t
transform of a time-dependent response function χ(t, t0 ) relating changes in z at time t to
forces at other times t0 : Z ∞
δz(t) = dt0 χ(t, t0 )F(t0 ). (3.30)
∞
Linear response functions are one of the most important tools for studying properties and
excitations of materials from metals to polymers. Interestingly, the response functions for
these much more complex functions not only share the same underlying mathematics as
our simple example but also exhibit the same behavior with narrow peaks in frequency
near elementary excitations.
3.1.3 Dissipation
It is the periodic forcing that sustains an oscillator’s motion and keeps it from succumbing
to the damping force. We need to calculate how much work is done by the external force in
one cycle. This energy is in the end lost to friction, which heats the system. Physically, you
can think of this as the energy dissipated by the oscillator. The rate at which the external
force does work is
dW
= F(t)3(t). (3.31)
dt
Recall that in the steady state, both F(t) and 3(t) oscillate at the same frequency ω and
are thus periodic functions of time with period T = 2π/ω. This allows us to calculate the
average power dissipated in a cycle:
1 T 1 T dF
Z Z
P̄d = dtF(t)ż(t) = − dt z(t), (3.32)
T 0 T 0 dt
P̄d (ω )
P̄d ∝ Q
15
Q = 16
10
1
2 P̄d ∆ω ≈ γ ∼ Q−1
0 ω
0 ω0
t
Fig. 3.5 Resonant power absorption as a function of frequency for three different values of Q.
The maximum Pmax and half maximum 12 Pmax as well as full width half maximum
frequency width ∆ω are indicated for the Q = 16 curve.
62 Resonance
t
where to obtain the second equation, we integrated by parts using the fact that F(T )z(T ) =
F(0)z(0). Using Eq. (3.15a) for z(t) and setting A0 (ω) = F0 |χ(ω)|, we obtain
F02
Z
P̄d = dt ω|χ(ω)| cos ωt sin[ωt − φ(ω)]
T
1 1
= ωF02 |χ(ω| sin φ(ω) = F02 ωχ00 (ω). (3.33)
2 2
It is a straightforward exercise to show that this result is identical the rate that the system
does work against dissipative forces. [See exercise 3.1.3]
Thus, we arrive at the very important result that the rate of energy dissipation is pro-
portional is proportional to ωχ00 (ω). For this reason, χ00 (ω) is sometimes referred to as the
dissipation. Note that χ00 (ω) is odd in ω so that ωχ00 (ω) is even. In thermodynamic equi-
librium, the power dissipation must be positive (energy is dissipated not created by friction
forces), implying that ωχ00 (ω) is positive. The positivity of ωχ00 (ω) in the present case is
associated with the positive sign of γ. Its sign was chosen so that the viscous or friction
force opposes motion of the oscillator mass. The sign is consistent with the transfer of en-
ergy to incoherent degrees of freedom of the fluid or medium responsible for friction and
to positive power absorption.
Figure 3.5 shows a plot of P̄d (ω) for three different values of Q. In contrast to the
frequency-dependent amplitude A0 (ω), the power absorption function P̄d (ω) is peaked ex-
actly at ω0 , and the full width at half maximum is precisely ∆ω = γ = ω0 /Q [Prob. 3.5] for
all γ, even when γ > 2ω0 and the solutions to the homogenous equation of motion has no
oscillatory part, though the peak can be quite asymmetric about ω0 . Setting ω = ω0 reveals
that the height of the power resonance curve is P̄d = (F02 /2mγ) ∝ Q.
Exercise 3.1.2 Using the fact that the force exerted by the oscillating weight on the
medium producing friction is −Fd and, thus, the work it does against friction is Γ3dx,
show that average power loss due to friction is equal to the average power, Eq. (3.33),
power spent by the the external force.
3.1.4 Energy stored in a driven oscillator
We noted at the outset of this chapter that a well-designed mechanical clock only needs a
very small kick each cycle to keep it going. What this means, of course, is that the energy
stored in the swinging pendulum of a mechanical clock is much greater than the energy
dissipated each cycle. Having already looked at the energy dissipated in a driven oscillator,
let’s calculated the energy stored.
The energy of a harmonic oscillator is the sum of its kinetic and potential energies, which
can be expressed in the following generic form
1 2 1
E= m3 + mω20 x2 , (3.34)
2 2
where the spring constant is given by k = mω20 . Using Eq. (3.15), z = A0 cos(ωt − φ) and
t
ż = −ωA0 sin(ωt − φ), so that Eq. (3.34) becomes
1
E = mA20 ω2 sin2 (ωt − φ) + ω20 cos2 (ωt − φ) .
h i
(3.35)
2
We are less interested in the instantaneous energy than we are in the average energy. There-
fore we average Eq. (3.35) over one cycle. The only time varying parts are the sine and
cosine terms. Over a single period, sin2 (ωt − φ) and cos2 (ωt − φ) oscillate between 0 and
1 with an average value of 21 . Performing these averages, Eq. (3.35) leads to the following
expression for the average steady state energy stored in a driven oscillator
1
Ē s = ml2 A20 ω2 + ω20 .

(3.36)
4
Using the expression for A0 given by Eq. (3.12), this becomes
F02 ω2 + ω20
Ē s (ω) = . (3.37)
4m ω2 − ω2 2 + γ2 ω2

0
To make direct comparison with P̄d (ω), we integrate Eq. (3.33) over one period to obtain
the energy dissipated in one cycle
2π
Z
Ēd (ω) = P̄d dt = T P̄d = P̄d (ω) . (3.38)
cycle ω
The ratio of the energy stored to the energy dissipated is given by
Es 1 ω2 + ω20
= , (3.39)
Ed 4π γω
which diverges for ω ω20 /γ and for ω ω0 and has a broad minimum at ω = ω0 ,
meaning that a little more energy is dissipated relative to the energy stored near resonance
than at other nearby frequencies. At ω = ω0 , the ratio Ē s /Ēd reduces to
Ē s (ω0 ) Q
= . (3.40)
Ēd (ω0 ) 2π
This is what we intuitively expect: higher Q oscillators dissipate a smaller fraction of their
energy than do lower Q oscillators. Equation (3.40) is often used as an alternative to Eq.
(2.28) as a definition of Q. It’s worth remembering that Q is proportional the energy stored
divided by the energy dissipated in a driven oscillator at resonance.
3.1.5 Transients
Suppose that an oscillator is initially at rest at its equilibrium position at z = 0 when a

periodic force is suddenly applied, starting at time t = 0. Because of its finite inertia, the
oscillator will not immediately start moving at its steady state velocity. Instead, it will take
some time for it to build up its motion. As the oscillator begins to move, it will exhibit some
transient behavior that over time will give way to the steady state oscillations we studied
in the previous sections. From our study of the damping of free oscillations of oscillators
in §2.1, we would expect any transients to die out after approximately Q oscillations, but
64 Resonance
t
as Q can easily be on the order of 104 , it could be quite awhile before a oscillator reaches
its steady state, that is, until its behavior is independent of the initial conditions. Here we
take a look at the transient motion of an oscillator as it builds up to its steady state.
To determine the transient behavior of the oscillator, we once again make use of the
principle of superposition, but in a different guise this time. Consider two solutions, z1 (t)
and z2 (t), to the motion of an oscillator: z1 (t) is a solution to Eq. (2.4), the equation of
motion for an unforced free oscillator, and z2 (t) is a solution to Eq. (3.4), the equation of
motion for a forced oscillator:
z̈1 + γ ż1 + ω20 z1 = 0 (3.41)

z̈2 + γ ż2 + ω20 z2 = (F0 /m) cos ωt . (3.42)
Suppose we sum these two solutions to create a new function z3 (t) = z1 (t) + z2 (t). Now,
let’s substitute z3 (t) into the left hand side of Eq. (3.42) and see what we get:
z̈3 + γ ż3 + ω20 z3 = z̈1 + γ ż1 + ω20 z1 + z̈2 + γ ż2 + ω20 z2

(3.43)
According to Eq. (3.41), the term inside of the first set of parentheses on the right hand
side of Eq. (3.43) is zero, while according to Eq. (3.42), the term inside the second set
of parentheses is F0 cos ωt. Therefore z3 (t) is also a solution to Eq. (3.4), the equation of
motion for a forced oscillator. This illustrates an important mathematical result: the sum of
a solution to the homogeneous differential equation and a of a solution of the inhomoge-
neous differential equation is also a solution of the inhomogeneous differential equation. In
particular, we can add the transient solutions we found in §2.1 to the steady state solutions
we found in §3.1.1 to obtain solutions to the full time-dependence of an oscillator.
This brings us to our method for finding the complete solutions: add the steady state
solution, which has no free integration constants, to the general solution for the undriven
damped oscillator, which has two undetermined integration constants. The two undeter-
mined integration constants are then chosen so that they satisfy the initial conditions, i.e.
the initial position and velocity of the oscillator. You might worry that in doing this, we’ve
missed something, that there is some other term we may have left out. Well, you needn’t
worry. There is a very important theorem concerning linear differential equations that says
that if you have found a solution to a differential equation that satisfies all the initial con-
ditions, then you have you have found the solution. It’s called the Uniqueness Theorem. It
is discussed in nearly every introductory text on differential equations.
As an example of how this works, consider the case where a sinusoidal forcing term is
applied starting at t = 0 to a oscillator that is initially at equilibrium: z(0) = 0 and ż(0) = 0.
Let’s consider the underdamped case where Q > 12 , so we will use the solutions from
§2.1.1. Adding together steady state and the transient general solutions, Eqs. (3.24) and
(2.23), respectively, with, as yet, undetermined integration constants A1 and A2 , gives the
general solution for the driven underdamped case:
z(t) = AR cos ωt + AI sin ωt + e−γt/2 (A1 cos ωu t + A2 sin ωu t) , (3.44)
where ωu is given by Eq. (2.16) and ω is the drive frequency. The amplitudes AR and AI
from the steady state solution can be written in terms of the real and imaginary parts of
t
the response function χ(ω) using Eqs. (3.25) and (3.26): AR = F0 χ0 (ω) and AI = F0 χ00 (ω).
Applying the initial conditions that z(0) = 0 and ż(0) = 0 to Eq. (3.44) yields
z(0) = F0 χ0 (ω) + A1 = 0 (3.45)
γ
ż(0) = ωF0 χ00 (ω) − A1 + ωu A2 = 0 . (3.46)
2
Solving for A1 and A2 gives
A1 = −F0 χ0 (ω) (3.47)
ω γ
A2 = − F0 χ00 (ω) − F0 χ0 (ω) . (3.48)
ωu 2ωu
Substituting A1 and A2 into Eq. (3.44) gives the solution for the specified initial conditions
γ
( " !#
z(t) = F0 χ (ω) cos ωt − e
0 −γt/2
cos ωu t + sin ωu t
2ωu
ω −γt/2
" #)
+ χ00 (ω) sin ωt − e sin ωu t . (3.49)
ωu
Figure 3.6 shows the solutions given by Eq. (3.49). The upper plot shows the results for
a driving frequency of ω = 0.9 ωu . Initially, both the (nearly) resonant frequency ωu and
the drive frequency ω are present in the response. Because they are nearly equal, the two
frequencies produce beats with a period of T B = 1/∆ν where ∆ν = |νR −ν| = |ωu −ω|/2π. As
time passes, the transient response at frequency ωu dies out, owing to the e−γt/2 damping
terms in Eq. (3.49), and the beats fade away. The damping occurs on a time scale of 2/γ or
equivalently after about Q = ω0 /γ oscillations (here, Q = 64).
The lower plot in Fig. 3.6 shows the results for a driving frequency at resonance ω = ωu .
In this case there are no beats because the oscillator is driven at the resonant frequency. In
this case, the steady state response builds up monotonically on a time scale of 2/γ (or ∼ Q
oscillations) towards the steady state solution.
Before moving on to applications of resonators, we pause briefly to reflect on the fact
that only the transient solution depends on the initial conditions; the steady state solution
is completely independent of the initial conditions. This means that no matter what the
starting conditions are, if a system is driven at a particular frequency and forcing, it will
end up in the same state. This is a property of linear oscillators, one that can be dramatically
violated for nonlinear oscillators.
3.1.6 Indirect forcing
Thus far we have considered the case that the periodic forcing is applied directly to an
oscillating mass. In many systems the forcing is applied indirectly, for example, by causing
the point from which a spring or pendulum is suspended to oscillate back and forth or up
and down. Resonance is still observed in such systems, but the asymptotic low and high
frequency behavior can be quite different from the case of direct forcing.
As an example of indirect forcing, let’s consider a pendulum, illustrated in Fig. 3.7,
whose pivot point is forced to oscillate horizontally back and forth according to z p =
z0 cos ωt. The pendulum is free to rotate about the pivot point. The mass at the end of
66 Resonance
t
z(t)
3
2
1
0 t/T
10 20 30 40 50 60 70 80
−1
−2
−3
z(t)
15
10
0 t/T
10 20 30 40 50 60 70 80
−5
−10
−15
t
Fig. 3.6 Transient response for an oscillator starting from rest (z = 0 and ż = 0) and a quality
factor Q = 64. (a) The driving frequency is ω = 0.9 ωu , slightly off resonance, and
beats appear at the difference frequency ∆ν = |νR − ν| = |ωu − ω|/2π. (b) The driving
frequency is ω = ωu , exactly at resonance. The response is much larger for the drive
frequency at resonance (note the difference in the z scale for the two plots). The time
axis in both cases is expressed in units of the natural period T = 2π/ω0 of the
oscillator.
the pendulum moves a distance s = lθ along an arc like a normal pendulum. For small
amplitude oscillations the motion of the mass at the end of the pendulum is to a good
approximation parallel to the displacement of the pivot point, so that the absolute position
of the mass is given by z p + s. The damping force on the pendulum is taken to be friction
from the hinge at the pivot point, and is proportional to ṡ, the velocity of the mass relative
to the pivot point (not the fixed reference frame). In this case, the equation of motion is
given by
d2 (z p + s) ds
m = −mgθ − Γ (3.50)
dt2 dt
Rearranging terms and using s = lθ, we obtain,
1 d2 z p
θ̈ + γ θ̇ + ω20 θ = − , (3.51)
l dt2
where ω0 = g/l and γ = Γ/m. Equation (3.51) bears a striking similarity to Eq. (3.4), the
p
equation of motion for the forced damped oscillator, but with the forcing term replaced by
t
−z̈ p /l, which is proportional to the acceleration of the pivot. Such a device is sometimes
used as a seismometer, a possibility we explore further in §3.1.6.
Response to sinusoidal acceleration

Equation (3.51) is valid for arbitrary accelerations z̈ p , so long as θ remains sufficiently
small. We begin our analysis by taking the acceleration to be sinusoidal. Later, we will
consider sudden impulses. We write the sinusoidal displacement as the real part of a com-
plex exponential
z p = z0 cos ωt = z0 Re{e−iωt } , (3.52)
and, as previously, will keep only the real part of any solution we obtain. Proceeding as
before, we assume a steady-state solution of the form θ(t) = Ae−iωt , where, following the
same reasoning, the frequency of the steady state response is the same as that of the driving
force. Substituting our trial solution θ(t) = Ae−iωt and Eq. (3.52) into Eq. (3.51) gives
z0 2
(−ω2 − iγω + ω20 )A = ω , (3.53)
l
where we have canceled the common factors of e−iωt . Solving for A, we obtain
(z0 /l) ω2
A= . (3.54)
ω20 − ω2 − iγω
Writing A = A0 eiφ , where a0 and φ are taken to be real, our solution takes the form
θ(t) = A0 e−i(ωt−φ) , (3.55)
t
Fig. 3.7 Pendulum with a moving pivot point. Gray image shows the equilibrium position of
pendulum and the pivot point.
68 Resonance
t
where
(z0 /l) ω2
A0 = q (3.56)
(ω20 − ω2 )2 + ω2 γ2
 γω 
 
φ = tan  2
−1 
 . (3.57)
ω0 − ω2
Figure 3.8 shows plots of Eqs. (3.56) and (3.57) for several values of Q = ω0 /γ. The res-
onance curves look similar to those we found previously for the forced damped oscillator.
However, the asymptotic behaviors here are different: for ω ω0 , the amplitude of the
oscillation goes to zero, i.e. A0 → 0, while for ω ω0 , the amplitude goes to a finite
value, i.e. A0 → z0 /l. The behavior at low frequency is consistent with what you might
have anticipated, namely, if the pendulum is moved back and forth too slowly, the pendu-
lum doesn’t oscillate but moves as a whole with its support: θ ≈ 0. When the pendulum
framework is moved back and forth very quickly at high frequencies, the mass on the end
of the pendulum basically remains (very nearly) still. Its angular amplitude A0 is just the
distance z0 the pivot point moves divided by the length of the pendulum. As a result, the
“oscillations" in the reference frame of the pivot point are exactly out of phase with the
pivot point.
Our quantitative analysis of the problem has shown that “too slowly" means at a fre-
quency much less than the natural oscillation frequency ω0 of the pendulum. Of course,
we might have guessed this from the outset, as there are only two time scales in the prob-
lem, ω−10 and γ , and it makes sense physically. If γ ω0 , meaning the damping is
−1
weak relative to the oscillations, the pendulum oscillates with respect to the pivot points.
If γ ω0 , meaning the damping is strong relative to the oscillations, the pendulum just
moves with the pivot point.
3.2 Resonant filters and detectors
Oscillators are frequently used as detectors that selectively filter out all signals except the
one of interest. One of the most common of all such detectors is the radio receiver, which
uses an electronic oscillator whose resonant frequency can be tuned to detect the one signal
that interests us among the many that are available. We discuss the radio receiver in §3.2.1.
Resonators as filters
Figure 3.9(a) shows that the amplitude of the response of the pendulum seismometer is
constant—or flat—above the resonance frequency ωR but falls off dramatically as the fre-
quency is reduced below ωR . The pendulum, in addition to serving as a seismometer, also
filters the input signals, responding to those with frequencies greater than the resonant fre-
quency while suppressing signals with frequencies less than the resonant frequency. Such
69 Resonant filters and detectors
t
A 0 (ω )
15
Q = 16
10
∆ω ∼ Q−1
0 ω /ω0
0 1 2 3
φ (ω )
π
Q = 16
π /2 Q = 1/2
0 ω /ω0
0 1 2 3
t
Fig. 3.8 Amplitude (top) and phase of pendulum response θ(t) to periodic shaking of the
pendulum for values of the quality factor Q from the darkest to the lightest curves: 12 ,
1, 2, 4, 8, 16. The dashed curve corresponds to Q = 1. The height of the peak of the
resonance curve (top) is given by A0 (ω0 ) = Q( f0 /ω20 ) and width ∆ω is proportional to
1/Q in the limit that Q 1. Phase (bottom) of the response of the pendulum. The
curves corresponding to Q = 1 are shown as dashed lines.
a filter is called a high-pass filter because it passes through the filter only those signals at
frequencies higher than the resonant frequency.
By contrast, the driven damped pendulum that we considered in §3.1.1 acts as a low pass
filter. Its behavior can be seen in Fig. 3.2 and in Fig. 3.9(b) above where the amplitude of
the response is shown for smaller values of Q.
We can make another kind of filter—a resonant filter—by using either one of these two
kinds of pendulums in a mode where we have Q 1. In this limit, the pendulums respond
only over a narrow range of frequencies centered around the resonant frequency.
All three kinds of filters, low-pass, high-pass, and resonant filters find wide applica-
tion in science and technology. In this section, we have studied devices that are in effect
mechanical filters. These same three kinds of filters are used in vastly different contexts,
however. There are electronic filters, optical filters, and acoustic filters. In fact, you can
find filters in almost any context in which you find oscillations. In the next section, we take
a look at radio receivers, which use resonant electronic filters. We will study the simplest
example, but one that is sufficient to illustrate the basic functioning filters in radios.
70 Resonance
t
(a) (b)
Q=2 Q=2
100 100
1 1
10−1 Q= 8 10−1 Q= 8
10−2 10−2
10−1 100 101ω /ωR 10−1 100 101ω /ωR
t
Fig. 3.9 (a) Amplitude vs. frequency for seismometer pendulum. (b) Amplitude vs. frequency
for driven pendulum described in §3.1.1 and by Eq. (3.12). For both graphs, the
values of Q plotted proceed from uppermost to lowermost trace: Q = 2, 1, 21 , 14 , and 18 .
3.2.1 Radio receiver
Radio signals permeate virtually every corner of the world we live in. At any given point in
space, the electric and magnetic fields associated with radio waves, coming from sources
near and far, fluctuate in time in a seemingly random manner, as illustrated in Figure
3.10(a). Some of the sources generate radio waves that, for lack of a better term, we call
noise. But many of those sources broadcast radio waves that, for lack of a better term,
E(t)
(a)
50
0 t (µ s)
0.2 0.4 0.6 0.8 1.0
−50
E(t)
(b) 5
0 t (µ s)
0.2 0.4 0.6 0.8 1.0
−5
t
Fig. 3.10 (a) Time dependence of the electric field of a radio signal at a particular point in
space. (b) Filtered radio signal. Note the difference in the scales of the electric field
E(t), and especially that the signal in part (b) is only a small fraction of the overall
signal shown in part (a).
t
we call information: music, news, and talk from commercial, nonprofit, and other sources.
The remarkable thing about radio and television receivers is that they are able to extract
from the morass that part of the signal that originates from a single source. A receiver does
this by filtering out essentially all the signals from sources not of interest. Figure 3.10(a)
shows what the time dependence of the electric field at a particular point in New York
City might look like; here we have included signals from more than 60 AM and FM ra-
dio stations. The modulation of the amplitude of the signal comes about because of the
interference between radio signals from the the various sources, each of which has its own
characteristic frequency. Because many of the broadcast frequencies are very near each
other, summing them produces the characteristic pattern of beats that are visible in the sig-
nal. Figure 3.10(b) shows that part of the same signal that includes only the signal coming
from a particular radio station, one that broadcasts at 99.5 MHz.2 It is the task of the radio
receiver to extract only that frequency from complex waveform shown in Figure 3.10(a).
The detection of radio signals is a complex process that depends on the size, shape, and
orientation of the antenna. We are not going to focus on that part of the detection of radio
waves. Suffice it to say that electric field of the the radio waves causes the electrons in the
antenna to oscillate back and forth producing a current whose time dependence follows the
electric field and thus resembles the signal shown in Figure 3.10(a). Our focus here is going
to be on how the electronics in a receiver selects the desired signal and filters out everything
else. The key, of course, is that each radio station broadcasts at a different frequency, so the
task of the receiver is to detect only the desired frequency and filter out the others. Such a
task is perfect for the electronic equivalent of the resonant filters introduced in the previous
sections.
At the heart of a radio receiver is a resonant filter, which in its simplest form consists of
an inductor, a capacitor, and a resistor connected to an alternating current or voltage source.
In the case of a radio receiver, it is the current, produced by radio waves, from the antenna
that acts as a current source. The current generated by radio waves is tiny in general and
thus must be amplified by a considerable amount to be processed and rendered audible.
In addition, a receiver must extract the audio signal from the radio wave, which it does
through a process called demodulation, which works differently for AM and FM signals,
for example. Thus, a radio receiver must detect radio waves, selectively filter, amplify, and
demodulate (and amplify again!) in order for us to hear the music. These four processes
aren’t always entirely independent of each other but for the moment, we shall focus only
on the filtering of a radio signal such that the signal from only one radio station is detected.
As noted above, different radio stations broadcast at different frequencies. The commer-
cial AM band corresponds to radio frequencies between about 520 kHz and 1.6 MHz, with
different radio stations’ signals typically separated by at least 10 kHz. The commercial
FM band corresponds to radio frequencies between about 87 MHz and 108 MHz, with
different radio stations’ signals typically separated by at least 100 kHz. So we will want
to design an electrical resonant filter that can tune to these frequencies with a sufficiently
sharp resonance such that nearby radio signals are filtered out.
2 99.5 MHz corresponds to a radio station in the FM or frequency modulated band. The audio signal is encoded
in small variations in the broadcast frequency. This means that the instantaneous broadcast frequency varies
by up to 20 kHz from the center broadcast frequency of 99.5 Mz, or about 0.02%.
72 Resonance
t
t
Fig. 3.11 Resonant circuit with a resistor, inductor, and capacitor in series.
To begin, we consider the circuit shown in Fig. 3.11. The voltage source serves as our
“signal," which, for the moment, we will take to have a sinusoidal time dependence:
V(t) = V0 cos ωt . (3.58)
The voltage drop across the capacitor is simply the instantaneous charge q on the capacitor
divided by its capacitance C. Similarly, the voltage drop across the inductor is L dI/dt and
across the resistor is IR. Thus we have
q dI
V(t) = + L + IR = V0 cos ωt. (3.59)
C dt
The current is related to the charge on the capacitor by I = dq/dt. Therefore, to obtain a
differential equation we can readily analyze, we take the derivation of Eq. (3.59) and obtain
d2 I dI 1
L + R + I = −ωV0 sin ωt . (3.60)
dt2 dt C
This has the same form as previous equations we have analyzed involving resonance, with
the current I as the dependent variable, L appearing in the place of mass, R in the place of
damping, and with C −1 providing the restoring term. Dividing through by L we obtain
d2 I R dI 1 ωV0
2
+ + I=− sin ωt . (3.61)
dt L dt LC L
Comparing this to Eq. (3.2), rewritten below for convenience
d2 z dz F0
+ γ + ω20 z = cos ωt , (3.62)
dt2 dt m
we can identify the damping term as γ = R/L, the resonance frequency as ω0 = (LC)−1/2 ,
which gives a quality factor of Q = ω0 /γ = [(L/R)/(RC)]1/2 = (L/C)1/2 /R, which is the
ratio of two time constants, L/R and RC.
The steady state response of the current I must oscillate at the driving frequency, just
as we found previously for mechanical oscillations. In electronics, it is conventional to
assume that the current has the form I(t) = I0 eiωt instead of I(t) = I0 e−iωt , so we follow
that convention here. Writing sin ωt = −ieiωt and substituting into Eq. (3.60), we obtain
1
−ω2 L I0 eiωt + iωR I0 eiωt + I0 eiωt = ωV0 ieiωt . (3.63)
C
t
Dividing both sides by iωI0 eiωt and rearranging terms a bit yields
V0 1
= R + iωL + . (3.64)
I0 iωC
The terms on the right hand side have the units of resistance. Taken together, they constitute
what is commonly called the complex impedence of the circuit, which is often denoted as
Z. Factoring out R, we can write the complex impedance as
V0 L 1
" #
Z= = R 1 + iω + (3.65)
I0 R iωRC
ω ω0
" !#
= R 1 + iQ − . (3.66)
ω0 ω
Thus, the current is given by
V0 iωt
I(t) = I0 eiωt = e (3.67)
Z
V0
= eiωt . (3.68)
R + iωL + 1/iωC
We can write I0 = V0 /Z in polar form I0 = |I0 |eiφ , where
 !2 −1/2
|V0 |  2 ω ω0 
|I0 | = 1 + Q − (3.69)
R  ω0 ω 

ω0 ω
" !#
φ = tan Q−1
− . (3.70)
ω ω0
Note that the current goes to zero at zero frequency and to a finite value for ω ω0 . Thus,
for the overdamped case when Q ≤ 1/2, the circuit would act as a high-pass current filter.
74 Resonance
t
Problems
3.1 A block of mass m = 100 g is hung from a spring whose force constant is k = 60 N/m.
The damping force on the block is given by Fd = −Γ3, where Γ = 3.0 N-m/s. The
mass is subject to an oscillatory driving force of F = F0 cos ωt, where F0 = 1.5 N.
(a) Write down the differential equation describing the steady state motion of the
block and solve for its amplitude and phase as a function of angular frequency.
Express your answers entirely in terms of symbols (no numbers).
(b) What are the resonant frequency and Q of the system?
(c) What is the amplitude of oscillation at the resonant frequency?
(d) At what frequencies is the amplitude of oscillation half of its maximum value?
(e) What are the limiting amplitudes of oscillation for ω much less than and much
greater than resonant frequency?
(f) What is phase difference between the driving force and the displacement at res-
onance? for ω = ω0 /2? for ω = 2ω0 ? Give your answers in both radians and
degrees and find the time delay between the peak in the driving force and the
peak in the mass displacement.
3.2 A 0.30-kg mass hangs at equilibrium from a spring suspended from a rod. Pulling the
mass down 4.0 cm from its equilibrium position requires a force of 5.0 N to stretch
the spring.
(a) Letting go of the mass after displacing it 4.0 cm from its equilibrium position,
you observe that its amplitude decays to 0.50 cm after 10 oscillations. What is
the Q of this system?
(b) Suppose that the system is now driven with a oscillating driving force of F(t) =
F0 sin ωt starting at time t = 0, where F0 = 0.50 N and the driving frequency
f = 2.5 Hz. The initial displacement of the mass is, once again, 4 cm below its
equilibrium position.
(i) Find an equation for the subsequent motion of the mass and then plot it as
a function of time.
(ii) What is the steady state amplitude the oscillations?
(iii) What is the initial period of the oscillations?
(iv) What is the ultimate period of the oscillations?
(v) How much energy is stored in the oscillating system at steady state?
(vi) How much energy is dissipated per period by the system at steady state?
3.3 Imagine an experiment similar to Millikan’s famous oil drop experiment in which
an oil droplet of mass m is suspended between two flat metal plates a distance h
apart and connected to a power supply that produces an oscillating voltage V(t) =
V0 cos ωt across the plates. Further assume that that the oil droplet carries a charge of
q, such that it experiences an oscillating force F = (qV0 /h) cos ωt from the oscillating
electric field. The damping force from the droplet moving through air is given by
75 Problems
t
Fd = −Γ3, where Γ is the damping constant and 3 is the droplet velocity. You may
ignore gravity in this problem. In this problem there is no spring-like restoring force.
(a) Write down the equation of motion for this system, including the oscillatory
driving force and proceed to solve it by assuming a trial solution of the form y =
Ae−iωt . Solving the resulting equation, show that the steady-state displacement
of the droplet is given by z(t) = A0 (ω) cos(ωt − φ) where
qV0
A0 = ,
m h ω ω2 + γ 2
p
γ
φ = tan−1 ,
−ω
and γ = Γ/m.
(b) The damping constant can be estimated from the equation Γ ≈ 3πηd, where
η ≈ 1.8 × 10−5 Pa-s is the viscosity of air (in SI units) and d is the diameter
of the droplet. Assuming a droplet diameter of d = 1 mm, a plate separation of
h = 1 cm, and an oil mass density of 0.8 g/cm3 , find the amplitude of oscillation
for a drive frequency of ω that is equal to the damping rate Γ/m and a voltage
amplitude of V0 = 100 V. Comment on the feasibility of such an experiment
given the numbers you find.
3.4 Consider an spring-and-mass oscillator such as that of Fig. 3.1, the top of whose
spring is moved up and down with an amplitude s(t) = a cos ωt. Set up the equation
of motion for this system, and show that the steady-state solution for the displace-
ment z(t) for the weight at the end of the spring is
ω20
z(t) = a q cos(ωt − φ).
(ω2 − ω20 )2 + ω2 γ2
3.5 Show that the power dissipation can be expresses as

1 1/Q
P̄d = F02 .
2mω0 ω ω0 2 1

ω0 − ω + Q2
Use this result to show that P̄d [Eq. (3.33)] has its maximum at ω = ±ω0 and that the
full width at half maximum is ωfwhm = ω0 /Q.
3.6 Consider a driven damped oscillator consisting of a mass hanging from a spring like
the one depicted in Fig. 3.1. Initially the system hangs at its equilibrium position.
At time t = 0 a sinusoidal oscillating force is applied to the mass: F(t) = F0 sin ωt.
Following the method used in §3.1.5, find the subsequent motion of of the oscillator
for a system with Q = 64 and a driving frequency of ω = 0.9ωu . Take the period of
the oscillator to be T = 2π/ω0 = 1 s. Compare your result with the result obtained
in §3.1.5 for a cosine driving force at times t QT and for times t QT .
3.7 Repeat Problem 3.6 except instead of applying a sinusoidal force starting at t = 0,
solve the problem for the case where the point from which the (massless) spring
hangs starts oscillating up and down at t = 0 with a time-dependent vertical dis-
placement given by y p (t) = y0 sin ωt. Take the damping force to be Fd = −Γ ẏ. Take
76 Resonance
t
the period of the oscillator to be T = 2π/ω0 = 1 s. Compare your result with the
result obtained in Problem 3.6 at times t QT and for times t QT .
3.8 The seismometer described in §3.1.6 can detect seismic motion in the horizontal
direction. To detect motion in the vertical direction, we need an oscillator capable of
motion in the vertical direction. A simple mass hanging from a suitable spring would
work, assuming the system is constrained to move only in the vertical direction.
(a) Write down Newton’s second law for this system and obtain the equation of
motion assuming the Earth and oscillator support oscillates vertically according
to the equation yE (t) = a cos ωt, where a is the amplitude of the motion. You
should follow the same kind of procedure used in §3.1.6 to set up the swing-gate
seismometer by defining coordinates relative to a fixed reference frame (not the
Earth!).
(b) Find the steady state motion, including the frequency-dependent amplitude A0 (ω)
and phase φ(ω) relative to the driving force. Plot A0 (ω) and φ(ω) as a function
of frequency for Q = 14 , 21 , 1, 2, 4, 8 and 16.
(c) Find and plot the response of the seismometer to a sudden impulse at time t = 0
assuming the mass is initially at rest at its equilibrium position. Choose values
for the spring constant, damping constant, and mass so that the instrument is
sensitive to frequencies from 0.1 to 1 Hz.
(d) Does this seismometer function as a low pass or high pass filter? Briefly explain.
3.9 In this problem, we consider an electrical resonant circuit consisting of two resistors,
a capacitor, and an inductor, driven by a sinusoidal voltage source V(t) = V0 cos ωt
across the open terminals shown in the figure below.
t
Fig. 3.12 Resonant RLC circuit.
(a) Find the differential equations describing for the current I in each branch of the
circuit and then solve for the voltage across the capacitor as a function of time.
Hint: Assume, as we did in §3.2.1, that the current has a complex amplitude and
is proportional to eiωt in each branch of the circuit and apply Kirchhoff’s law for
the current at each junction.
(b) What are the resonant frequency and Q for this circuit in terms of R, L, and C?
(c) On a single graph, plot the amplitude and phase of the voltage across the capaci-
tor in the steady state as a function of ω/ω0 assuming a values of Q = 1, 10, and
100. Be sure to label which curves belong to which values of Q.
77 Problems
t
3.10 (French 4-3) An object of mass 0.2 kg is hung from a spring whose spring constant
is 80 N/m. The body is subjected to a resistive force given by −b3, where 3 is the
velocity (m/sec) and b = 4 N-m−1 s.
1 Set up the differential equation for free oscillations of the system and find the
period of such oscillations [Ans: T = 5 √π 3 = 0.36s]
2 The object is subjected to a sinusoidal force given by F(t) = F0 sin ωt, where
F0 = 2 N and ω = 30 s−1 . In the steady state, what is the amplitude of forced
oscillations? [Ans: A = 1.28 cm]
3.11 (French 4-11) Consider a damped oscillator with m = 0.2 kg, b = 4 N-m−1 -s and
k = 80 N/m. Suppose the oscillator is driven by a force F = F0 cos ωt, where F0 = 2
N and ω = 30 s−1 .
1 What are the values of A and δ of the steady-state response described by z(t) =
A cos(ωt − δ) [Ans: A = 1.3 cm, δ = 130◦ ]
2 What is the mean power input. [Ans: P = 18/61 = 0.295 J/s.]
3 How much energy is dissipated against the resistive force in one cycle? [Ans:
Energy/cycle = PT = 0.062 J]
4 Normal Modes
Up to this point, we have considered oscillations of systems with with only a single degree
of freedom—the angle θ a pendulum makes with the vertical, the distance between two
atoms in a diatomic molecule, etc. However, most things that oscillate have many degrees
of freedom and typically oscillate at many different frequencies. Consider a string on a
guitar or piano. It can vibrate at its fundamental frequency or its first, second, or higher
overtones. The same is true of every other kind of musical instrument, or in fact, almost
any other object that vibrates—a plate, a goblet, a bridge, or even a skyscraper. Typically,
for each frequency of vibration, there is a characteristic pattern of displacements of the
vibrating object. These patterns of displacements and the characteristic frequency that goes
with each of them are known as normal modes.
4.1 Systems with two degrees of freedom
In this chapter we begin our study of normal modes with examples of systems that have
only two (or three) degrees of freedom, as these serve to introduce the basic concepts. Once
you understand those systems, we will move on to systems with more degrees of freedom,
where the real power of normal mode analysis, both conceptually and mathematically,
is revealed. As might be expected, when the number of degrees of freedom increases,
the algebra can become very tedious. To keep the algebra manageable, we introduce the
the tools of linear algebra: matrices, vectors, and eigenvalues. The mathematics of linear
algebra also introduces a conceptual framework for better understanding the physics of
systems with many degrees of freedom. Don’t worry if you are not too familiar with linear
algebra. We will introduce the tools and concepts of linear algebra as we need them.
4.1.1 Two coupled pendulums
We begin our exploration of normal modes by considering two identical pendulums con-
nected by a spring, as shown in Fig. 4.1. The spring connecting the two masses is un-
stretched when the two pendulums hang vertically, that is, when θ1 = θ2 = 0. Gravity
provides a restoring force for the pendulums of F1g = −mg sin θ1 and F2g = −mg sin θ2 ,
respectively. When stretched or compressed, the spring exerts a force F1s ' k(x2 − x1 ) on
mass 1 and F2s ' −k(x2 − x1 ) on mass 2 in the limit that θ 1, where we can neglect the
effects of the vertical extension of the spring. The sign convention is chosen so that when
78
79 Systems with two degrees of freedom
t
t
Fig. 4.1 Two identical pendulums of length l and mass m connected by a spring with spring
constant k.
the spring is stretched, it exerts a force to the right on mass 1 and to the left on mass 2.
Summing the forces on each mass, we obtain the following equations of motion:
x1
F1 = −mg + k(x2 − x1 ) = m ẍ1 (4.1)
l
x2
F2 = −mg − k(x2 − x1 ) = m ẍ2 , (4.2)
l
where we have used the approximations sin θ1 ' x1 /l and sin θ2 ' x2 /l, which are valid
when x1 , x2 l. Collecting terms together and rearranging gives
mg
m ẍ1 = − + k x1 + kx2 (4.3)
l
mg
m ẍ2 = kx1 − + k x2 , (4.4)
l
Equations (4.3) and (4.4) are the equations of motion for two pendulums connected by
a spring. They are coupled equations in that the dynamical variables x1 and x2 appear in
both equations. Physically, this means that exciting pendulum 1, say by giving it some
finite displacement, x1 , 0, or finite velocity, 31 , 0, will in general cause pendulum 2 to
move so that x2 , 0 and 32 , 0. Thus we say that the dynamical variables x1 and x2 are
coupled. Mathematically, the fact that these equations are coupled means that we will need
to solve for both x1 and x2 at the same time.
We are going to solve Eqs. (4.3) and (4.4) in two stages. First we will find the natural
frequencies of oscillation, the so-called normal frequencies. This particular system has two
degrees of freedom, the coordinates x1 and x2 of the two pendulums. Because there are two
degrees of freedom, there will be two normal frequencies. In general, the number of normal
frequencies is equal to the number of degrees of freedom, although in some cases there can
be fewer normal frequencies.
Second, we will find equations of motion for x1 and x2 . We will do this using a some-
what indirect approach that benefits from the hindsight that comes from others having
solved these kinds of problems many times before. The approach consists of defining new
dynamical variables qα and qβ that are linear combinations of x1 and x2 . By choosing those
linear combinations wisely, we will be able to turn the two coupled differential equations,
80 Normal Modes
t
Eqs. (4.3) and (4.4), into two independent uncoupled equations that are in the precise form
of the equation of motion for a single simple harmonic oscillator that we encountered in
Chapter 1 (see Eq. (1.5)). And we will be able to use the same sinusoidal solutions that we
developed in Chapter 1.
Normal frequencies
To solve Eqs. (4.3) and (4.4), we follow the same procedure we used for the case of a single
oscillator and assume complex sinusoidal solutions of the form
x1 = a1 e−iωt (4.5)
x2 = a2 e −iωt
. (4.6)
Substituting Eqs. (4.5) and (4.6) into Eqs. (4.1) and (4.2) and collecting terms gives
g k k
!
+ − ω2 a1 − a2 = 0 (4.7)
l m m
k g k
!
− a1 + + − ω2 a2 = 0 . (4.8)
m l m
These equations are satisfied only when the determinant of the matrix formed by the coef-
ficients of the amplitudes a1 and a2 is zero,1 that is, when
g/l + k/m − ω2

−k/m
=0 (4.9)

−k/m g/l + k/m − ω2

which gives
!2 !2
g k k
+ − ω2 − =0. (4.10)
l m m
Equation (4.10) is known as the secular or characteristic equation for the system of cou-
pled equations formed by Eqs. (4.7) and (4.8). Solving the secular equation2 for ω gives
q
ω = ± (ω2p + ω2s ) ± ω2s , (4.11)
√
where ω p = g/l and ω s = k/m. Equation (4.11) yields two frequencies, which we write
p
out explicitly
g
r
ωα = ω p = (4.12)
l
r
2 2 1/2 g 2k
ωβ = (ω p + 2ω s ) = + . (4.13)
l m
1 You are asked to prove this assertion, which is a special case of Cramer’s rule, in Exercise 4.1.1. It’s handy to
remember it! For an introduction and review of matrices and linear algebra, see Appendix A.
2 You can blindly expand Eq. (4.10) and solve it using the quadratic formula, or you can save yourself a lot of
work by noticing that it consists of two perfect squares, allowing you to simply write down the solution by
inspection, saving work and avoiding algebra errors.
t
These two frequencies ωα and ωβ are the normal frequencies of this system. To under-
stand their physical significance, we need to find the normal coordinates qα and qβ and the
normal modes of the system.
Exercise 4.1.1 Cramer’s Rule: Given a pair of coupled linear equations of the form
B11 a1 + B12 a2 = 0 (4.14)

B21 a1 + B22 a2 = 0 , (4.15)
show that the two equations are simultaneously satisfied when
B11 B22 − B12 B21 = 0 . (4.16)
If we define a matrix from the coefficients Bi j of Eqs. (4.14) and (4.15)
B11 B12
!
B= , (4.17)
B21 B22
then the determinant of the 2 × 2 matrix B is written as |B| or det B and is defined as
B11 B22 − B12 B21 . Thus, Eq. (4.16) can be written as |B | = det B = 0.
Normal coordinates
An interesting and extremely useful feature of systems of coupled linear oscillators is that
it is always possible to define new coordinates qα and qβ that are linear combinations of the
original coordinates x1 and x2 such that the equations of motion for these new dynamical
variables are uncoupled. Analyzing the composite system in terms of these new dynamical
variables results in an enormous mathematical and conceptual simplification of the system.
The new dynamical variables are sufficiently important that they are given a special name:
normal coordinates. Fancy words! Let’s illustrate what they mean by finding the normal
coordinates of the system of coupled pendulums we have been studying.
The idea is to recast Eqs. (4.1) and (4.2) so that we obtain two equations, each of which
involves a single dynamical variable that does not appear in the other equation. This is
readily accomplished by summing and subtracting Eqs. (4.1) and (4.2). Summing Eqs.
(4.1) and (4.2) gives
x1 + x2
−mg = m( ẍ1 + ẍ2 ) . (4.18)
l
Subtracting Eq. (4.1) from Eq. (4.2) gives
x2 − x1
−mg − 2k(x2 − x1 ) = m( ẍ2 − ẍ1 ) , (4.19)
l
Collecting and rearranging terms, these equations become
g
( ẍ1 + ẍ2 ) + (x1 + x2 ) = 0 (4.20)
!l
g 2k
( ẍ2 − ẍ1 ) + + (x2 − x1 ) = 0 . (4.21)
l m
82 Normal Modes
t
Defining the new dynamical variables
qα = x1 + x2 (4.22)
qβ = x2 − x1 , (4.23)
Eqs. (4.20) and (4.21) can be rewritten in more compact forms
g
q̈α + qα = 0 (4.24)
l!
g 2k
q̈β + + qβ = 0 . (4.25)
l m
The new coordinates qα and qβ are linear combinations of the original coordinates x1 and
x2 . Moreover, with this choice of qα and qβ , the equations of motion, Eqs. (4.22) and (4.23),
become decoupled: Eq. (4.22) is the equation of motion for qα only, and Eq. (4.23) is the
equation of motion for qβ only. This is important. It means that the dynamics of qα , that
is how qα develops in time, does not depend on qβ , and vice versa; qα (t) and qβ (t) are
completely independent of each other. The new dynamical variables qα and qβ are called
the normal coordinates of the system.
The equations of motion Eqs. (4.22) and (4.23) for the normal coordinates qα and qβ are
the same as the equation of motion for a simple harmonic oscillator (cf. Eq. (1.5)),
q̈α + ω2α qα = 0 (4.26)

q̈β + ω2β qβ =0, (4.27)
where ωα and ωβ are constants—the oscillation frequencies—of the different normal modes.
This is a general feature of normal coordinates: their corresponding equations of motion are
always the equation of motion of the simple harmonic oscillator. The only distinguishing
feature is that each normal coordinate has its own characteristic frequency of oscillation.
Comparing Eqs. (4.22) and (4.23) with Eq. (1.5), we can immediately read off the oscil-
lation frequency corresponding to each normal coordinate: qα oscillates at the frequency
ωα = sqrtg/l; qβ oscillates at the frequency ωβ = g/l + 2k/m. Of course, these are just
p
the frequencies of the normal modes we identified in Eqs. (4.12) and (4.13). This is also
a general feature of normal coordinates: the natural oscillation frequency of each normal
coordinate corresponds to one of the normal frequencies of the system.
4.1.2 Normal modes of two coupled pendulums
The two frequencies ωα and ωβ and the two normal coordinates qα and qβ correspond to
two separate modes of oscillation of the system of two coupled pendulums. In the first
mode, both pendulums oscillate at a frequency ωpα ; in the second mode, both pendulums
oscillate at a frequency ωβ . Since ωα = g/l < g/l + 2k/m = ωβ , the α mode is slower
p
than the β oscillation mode. Let’s examine each of these modes individually.
t
The slow mode
We start by examining the slow mode, which oscillates at the frequency ωα . The solutions
are given by summing the positive and negative frequency ±ωα solutions given by Eqs.
(4.5) and (4.6)
x1 (t) = a1 e−iωα t + a∗1 eiωα t (4.28)

−iωα t
x2 (t) = a2 e + a∗2 eiωα t , (4.29)
where we recall that the two terms in each sum must be complex conjugates of each p other
for the solutions x1 (t) and x2 (t) to be real (see §1.5.2). Substituting ω = ωα = g/l into
Eq. (4.7), we find that a1 = a2 for the α (slow) mode. With a1 = a2 , Eqs. (4.28) and (4.29)
are identical. That is, for the slow mode, x1 (t) = x2 (t).
The equation of motion for the slow mode is given by Eq. (4.24) or equivalently Eq.
(4.26). As the equation of motion is identical to that of a single simple harmonic oscillator,
so the general solution can be written as
qα (t) = 2A cos(ωα t − δ) , (4.30)
where we chosen to insert a factor of 2 into the definition of the amplitude A (see Eq.
(1.70) and §1.5.2). Since qα (t) = x1 (t) + x2 (t) and x1 (t) = x2 (t), the general solutions for
the equations of motion for the individual masses oscillating in the slow mode are
x1 (t) = x2 (t) = A cos(ωα t − δ) . (4.31)
According to Eq. (4.31), the two masses have exactly the same equation of motion, mean-
ing that they oscillate in phase with each other at a frequency ωα . Because they move in
phase with each other, the spring connecting them is neither stretched nor compressed; the
distance between the masses remains unchanged. Thus, gravity provides the only restoring
force and the frequency ωα = g/l of the coupled system is the same as for the uncoupled
p
system.
t
Fig. 4.2 Modes of two pendulums coupled by a spring.
84 Normal Modes
t
The fast mode
Following the same procedure we used for analyzing the slow mode, the solutions for the
fast mode are given by summing the positive and negative frequency ±ωβ solutions given
by Eqs. (4.5) and (4.6)
x1 (t) = a1 e−iωβ t + a∗1 eiωβ t (4.32)
−iωβ t
x2 (t) = a2 e + a∗2 eiωβ t . (4.33)
Substituting ωβ = g/l + 2k/m into Eq. (4.7), we find that a1 = −a2 for the β (fast) mode.
p
With a1 = −a2 , x1 (t) = −x2 (t).

The equation of motion for the fast mode is given by Eq. (4.25), or equivalently by Eq.
(4.27), and has the same general solution
qβ (t) = 2B cos(ωβ t − δ) . (4.34)
Since qβ (t) = x1 (t) + x2 (t) and x1 (t) = −x2 (t), the general solutions for the equations of
motion for the individual masses oscillating in the slow mode are
x1 (t) = −x2 (t) = B cos(ωβ t − δ) . (4.35)
For the fast mode, we see that the two pendulums oscillate synchronously exactly out of
phase with each other, as shown in the right panel in Fig. 4.2. In the fast mode, the spring
is alternately stretched and compressed each cycle. In this case, both the spring and gravity
provide the restoring force, which is reflected in the equation for the normal frequency
ωβ = (g/l + 2k/m)1/2 = (ω2p + 2ω2s )1/2 .
Solving coupled oscillator problems using normal modes

Because the equations of motion are linear, linear combinations of the solutions, including
those of different normal modes, are also solutions to the equations of motion. This is
the principle of linear superposition. In fact, all possible solutions of any coupled linear
oscillator can be expressed as a linear combination of the solutions to the individual normal
modes.
This brings us to the practical question of how best to find the solution to a coupled
oscillator problem for a particular set of initial conditions. In general, the simplest way to
solve such a problem is (i) first find the normal coordinates of the system, (ii) then express
the initial conditions in terms of those normal coordinates, (iii) and finally, find the solution
for each normal mode subject to the initial conditions for that normal coordinate. Once you
have the solutions in terms of the normal coordinates qα (t) and qβ (t), find the equations for
x1 (t) and x2 (t) by inverting the equations for qα (t) and qβ (t).
Let’s illustrate this procedure with an example. Before doing so, however, we note that
each normal mode is governed by an equation of motion for a simple harmonic oscillator:
q̈α + ω2α qα = 0. (4.36)
Because the equation of motion for each normal coordinate is the equation of motion of a
simple harmonic oscillator, the solution for each normal mode can be written in any one of
t
several forms. First there is the complex exponential form which is useful for setting up the
problem and for finding the normal frequencies and normal modes. But for constructing
useful solutions, it is usually better to use real functions. The two most common forms that
we have employed are
qα (t) = a1 cos ωα t + a2 sin ωα t (4.37)
qα (t) = A cos(ωα t − δ) , (4.38)
where α refers to a particular normal mode. In either case, there are two integration con-
stants, consistent with the second order equation of motion Eq. (4.26), to determine from
the initial conditions, the amplitudes a1 and a2 in the first case and the amplitude A and
phase δ in the second. Which form is used to solve a particular problem is to some extent
a matter of taste as the two descriptions are completely equivalent. However, at the final
stage of finding a particular solution to a normal mode problem with specific initial condi-
tions, it is generally simpler to use sines and cosines rather than the cosine with a phase.
We illustrate how this works by applying it to obtain solutions for two coupled pendulums.
Example 4.1 Find the equations of motion of two pendulums of mass m coupled by
a spring with a spring constant k for the initial conditions x1 (0) = 1, x2 (0) = 0, and
v1 (0) = v2 (0) = 0.
Solution
We solve this problem in terms of the normal coordinates, which we define using the un-
normalized form given by Eqs. (4.22) and (4.23), and using the sine and cosine form of the
solutions given by Eq. (4.38):
qα (t) = A1 cos ωα t + A2 sin ωα t (4.39)
qβ (t) = B1 cos ωβ t + B2 sin ωβ t . (4.40)
From Eqs. (4.22) and (4.23), the initial conditions for the normal coordinates are
qα (0) = A1 = x1 (0) + x2 (0) = 1 (4.41)
qβ (0) = B1 = x2 (0) − x1 (0) = −1 . (4.42)
Similarly, the initial conditions for the velocities gives
q̇α (0) = ωα A2 = ẋ1 (0) + ẋ2 (0) = 0 (4.43)
q̇β (0) = ωβ B2 = ẋ2 (0) − ẋ1 (0) = 0 , (4.44)
which means that A2 = B2 = 0. Therefore, the full solutions for the normal modes are:
qα (t) = cos ωα t (4.45)
qβ (t) = cos ωβ t . (4.46)
Using Eqs. (4.22) and (4.23), these solutions give
x1 (t) = 21 [qα (t) + qβ (t)] = 12 [cos ωα t + cos ωβ t] (4.47)
1 1
x2 (t) = 2 [qα (t) − qβ (t)] = 2 [cos ωα t − cos ωβ t] . (4.48)
86 Normal Modes
t
x1
1
0 t
5 10 15 20 25 30 35 40
−1
x2
1
0 t
5 10 15 20 25 30 35 40
t
−1
Fig. 4.3 Displacements of two pendulums coupled by a spring for the case where the natural
periods of the pendulums and spring are T p ≡ 2π/ω p = 3 s and T s ≡ 2π/ω s = 3 s,
respectively, which, from Eqs. (4.12) and (4.13), corresponds to T α = 2π/ωα = 3 s and
T β = 2π/ωpβ = 1.73 s. For the
√ numerical values chosen here,
T α /T β = 1 + (T p /T s )2 = 2. Because this ratio is an irrational number, the
displacements of the two pendulums never repeat.
These solutions can produce quite complicated trajectories for the two pendulums, as
shown, for example, in Fig. 4.3. For this case, we have chosen the natural frequencies ω p
and ω s (or periods T p and T s ) of the spring and pendulums to be comparable to each other.
In this limit the force exerted by the spring on the two mass is comparable to the force of
gravity, and we say that the two pendulums are strongly coupled. In the next section, we
examine the particularly interesting limit of weak coupling.
4.1.3 Weak coupling
It often turns out that the different degrees of freedom of a system are only weakly coupled
to each other. By this we mean that the characteristic energy (or force) with which different
degrees of freedom interact is small compared to the characteristic energies (or forces)
associated with the oscillations of those degrees of freedom. This arises, for example, for
the case of the two pendulums coupled by a spring when the characteristic energy E s of
compressing the spring is much smaller than the gravitational energy E p associated with
the pendulums’ rise. In the limit of small amplitude oscillations, these two energies can be
written as E s ∼ mω2s a2 and E p ∼ mω2p a2 (see Eq. (1.52)). The weak coupling limit thus
corresponds to cases for which E s E p , or equivalently, ω2s /ω2p 1.
The trajectories for two weakly coupled pendulums are shown in Fig. 4.4. The trajecto-
ries resemble the familiar pattern of beats observed when two sound waves of very nearly
t
x1
1
0 t
5 10 15 20 25 30 35 40
−1
x2
1
0 t
5 10 15 20 25 30 35 40
t
−1
Fig. 4.4 Displacements of two pendulums weakly coupled

by a spring. Gray
dashed lines
show modulated amplitudes: top, ± cos 12 ∆ω t and bottom, ± sin 21 ∆ω t .

the same frequencies interfere. Similarly, a beat pattern appears in the trajectories of cou-
pled two pendulums when the normal mode frequencies are very nearly the same. In the
present case, the difference in the normal mode frequencies is given by
∆ω = ωβ − ωα = (ω2p + 2ω2s )1/2 − ω p (4.49)

1/2
= ω p 1 + 2ω2s /ω2p −1 (4.50)
= ω p 1 + ω2s /ω2p − 21 ω4s /ω4p + ... − 1

h i
(4.51)
ωs
!
' ωs (4.52)
ωp
where we have used the binomial expansion Eq. (1.89) in the penultimate step and have
dropped the higher order terms in the expansion in the final step, which is a good approx-
imation in the limit that ω2s /ω2p 1. In that limit, it is also clear that ∆ω ωα , ωβ , and
therefore that the two normal mode frequencies ωα and ωβ are nearly the same. Defining
ω̄ = 12 (ωα + ωβ ), we can rewrite Eqs. (4.47) and (4.48) as
x1 (t) = 12 cos (ω̄ − 12 ∆ω)t + 12 cos (ω̄ + 21 ∆ω)t

h i h i
(4.53)
1

= cos 2 ∆ω t cos ω̄t (4.54)
x2 (t) = 21 cos (ω̄ − 12 ∆ω)t − 12 cos (ω̄ + 21 ∆ω)t
h i h i
(4.55)
= sin 21 ∆ω t sin ω̄t ,

(4.56)
where we have used the cosine angle addition formulas to obtain Eqs. (4.54) and (4.56).
Equations (4.54) and (4.56) show that x1 (t) and x2 (t) oscillate at the mean frequency ω̄ of
the two normal modes with a sinusoidally modulated amplitude with nodes spaced by a
time 2π/∆ω.
The trajectories plotted in Fig. 4.4 show an interesting pattern where the two oscillators
88 Normal Modes
t
“take turns” oscillating; when one oscillates with its maximum amplitude, the other is
essentially still. This is one of the hallmarks of two weakly coupled oscillators and can be
observed in a very wide variety of systems.
4.1.4 Energy in normal modes
From the previous discussion it should be clear that energy is passed back and forth be-
tween mass 1, mass 2, and the spring. Consider the case of the weakly coupled pendulums
discussed in the previous section. The gravitational potential energy and kinetic energy of
mass 2 are zero at t = 0 and again at t = 2π/∆ω, since the displacement and velocity of
mass 2 are zero at those times. On the other hand, the gravitational potential energy of
mass 1 is at a maximum, as it is maximally displaced from its equilibrium position at those
times. The potential energy of the spring is non-zero as it is stretched from its equilibrium
position. As time passes, all these energies change, even as the total energy of the system
must remain constant, since no energy is dissipated.
We can make an even stronger statement about energy conservation: the total energy
associated with each normal mode is constant in time. This follows from the fact that
normal modes are independent of each other. Because they are independent of each other,
energy is not passed back and forth between the modes but is independently conserved for
each mode. The amount of total energy stored in each normal mode depends on the initial
conditions, which determine the degree to which mode is excited.
Determining exactly how much energy is to be associated with each mode can be tricky.
For example, what is the kinetic energy of the α mode as a function of time? For a single
particle of mass m and velocity 3, the kinetic energy is simply 21 m32 . For the α mode, we
might take the velocity to be equal to q̇α (t), but what mass should we use? The mass m of a
single particle? or the mass 2m of two particles (since two masses are involved)? or some
other value? And what is the potential energy for mode 1? It seems clear that we should not
include the potential energy of the spring, since the α mode does not involve stretching the
spring. But how do we divide the gravitational potential energy of the two masses divided
between the α and β modes?
Let’s start by determining the kinetic energy Kα (t) as a function of time for the α mode.
We can do this by calculating the kinetic energies of particles 1 and 2 assuming that only
the α mode is excited initially, since what happens in the α mode is independent of what
happens in the β mode. For the spring-coupled pendulums, the displacements of the two
masses are given by Eqs. (4.47) and (4.48)
x1 (t) = 12 [qα (t) + qβ (t)] (4.57)

1
⇒ x1 (t) = 2 qα (t) (4.58)
1
x2 (t) = 2 [qα (t) − qβ (t)] (4.59)
1
⇒ x2 (t) = 2 qα (t) . (4.60)
The kinetic energy in the α mode is just the sum of the kinetic energies of masses 1 and 2
89 Matrix formulation of normal modes
t
under the assumption that the β mode is not excited. Thus, we obtain
Kα (t) = 21 m ẋ12 + 12 m ẋ12 (4.61)

2 2
= 21 m 12 q̇α + 12 m 12 q̇α

(4.62)
1 2
= 2 (m/2)q̇α . (4.63)
Based on our choice of normal coordinates, the effective mass meff associated with the qα
mode is m/2 according to Eq. (4.63).
The gravitational potential energy in the α mode is obtained by summing the potential
energy of particles 1 and 2 assuming the β mode is not excited:
Uα (t) = mgy1 + mgy2 (4.64)

= mgl(1 − cos θ1 ) + mgl(1 − cos θ2 ) (4.65)
= mgl 1 − 1 − 12 θ12 + ... + mgl 1 − 1 − 21 θ22 + ...
h i h i
(4.66)
' 12 mgl θ12 + θ22

(4.67)
mg 2
x1 + x22

' (4.68)
2l
mg 1 2 1 2
' qα + 2 qα (4.69)
2l 2
1 mg 2

= q , (4.70)
2 2l α
where we have used the small-amplitude approximation. Equation (4.70) is in the usual
form 21 keff x2 for the potential energy of a simple harmonic oscillator with an effective
spring constant keff = (m/2)(g/l) = meff ω2α , which is consistent with the usual relation for
an ideal spring simplep harmonic oscillator with an effective mass of m/2 and an oscillation
frequency of ωα = g/l. We note that our results for the effective mass meff and effective
spring constant keff depend on our choice of normal coordinates, which are determined
only to within an overall multiplicative constant.
Exercise 4.1.2 Following the same procedures used above, find the small angle expres-
sions for Kβ (t) and Uβ (t), the time-dependent kinetic and potential energies of the β
mode for the spring-coupled pendulums. In particular, show that the effective mass of
the β mode is given by m/2, and that the expression for the effective spring constant
is consistent with this value and the value of ωβ . Then show by explicit calculation
that the total energy for each of the modes Eα = Kα (t) + Uα (t) and Eβ = Kβ (t) + Uβ (t)
are constants independent of time.
4.2 Matrix formulation of normal modes
The coupled pendulum problem discussed in the previous sections is actually one of the
simplest examples of coupled oscillators. Because it is so simple, it is relatively easy to
90 Normal Modes
t
analyze. The normal coordinates qα and qβ turn out to be about the simplest set of linear
combinations of the particle coordinates x1 and x2 imaginable: qα = x1 +x2 and qβ = x1 −x2 .
So it was fairly easy to guess their form just by looking at the equations of motion, Eqs.
(4.3) and (4.4).
Unfortunately, finding the normal modes for most coupled oscillator problems is not so
simple. Moreover, the algebra associated with normal modes can get fairly tedious and
complicated. While tedious algebra is simply a fact of life sometimes, it can obscure the
physics and impede conceptual development of a subject. Introducing more compact math-
ematical notation, properly conceived, can improve matters. Consider, for example, the in-
troduction of vectors into physics. Using vector notation, we can write Newton’s second
law as
F = ma , (4.71)
which with one equation neatly summarizes three equations written in algebraic form:
F x = ma x Fy = may Fz = maz . (4.72)
Writing Newton’s second law in vector form can greatly simplify algebraic manipulations,
even if eventually we might need to use the component forms to do some of the mathemat-
ical manipulations.
There is another perhaps less apparent advantage to writing Newton’s second law in
vector form. It helps us to think about the physics of Newton’s second law in terms of
vectors—arrows—that point some direction in space without reference to any particular
coordinate system. Thus we don’t have to get bogged down in defining a coordinate system
and thinking only in components along one axis or another. We can think about the essential
elements of the problem geometrically, using helpful and compelling pictures. Enabling
us to think on this conceptual level is every bit as important as facilitating mathematical
manipulations.
For these reasons, both practical and conceptual, we embark below on a reformulation
of normal modes in terms of matrices and the mathematics of linear algebra.
4.2.1 Two equal-mass particles
Consider two equal masses m constrained to move along a frictionless surface as shown in
Fig. 4.5. The equations of motion for the two masses are
m ẍ1 = −ka x1 + kb (x2 − x1 ) (4.73)

m ẍ2 = −kb (x2 − x1 ) . (4.74)
t
Fig. 4.5 Horizontal masses on a frictionless surface coupled by two springs.
t
To determine the sign on the (x2 − x1 ) terms, note that when (x2 − x1 ) > 0, the kb spring
is stretched and it pushes mass 1 to the left in the negative direction and mass 2 to the
right in the positive direction. Collecting the x1 and x2 terms together, we rewrite these two
equations in standard form
m ẍ1 = −(ka + kb )x1 + kb x2 (4.75)
m ẍ2 = kb x1 − kb x2 . (4.76)
This equation can be written in compact matrix form with the introduction of the stiffness
matrix k and the column vector x
ka + kb −kb x1
! !
k= , x= . (4.77)
−kb kb x2
With these definitions, the two equations of motion can be expressed as a single matrix
equation
1
ẍ(t) = − k x(t) ≡ −D x(t) , (4.78)
m
where
(ka + kb )/m −kb /m ω2 + ω2 −ω2b
! !
D= = a 2 b (4.79)
−kb /m kb /m −ωb ω2b
is the dynamical matrix with ωa ≡ ka /m and ωb ≡ kb /m. Equation (4.78), ẍ(t) = −D x(t),
together with the definitions of D and x, is just a compact way of rewriting Eqs. (4.75) and
(4.76). It is very useful, however, because it permits us to use the mathematical machinery
of linear algebra to solve these equations.
We pause here to take note of an important mathematical property of the stiffness matrix
k : it is a symmetric matrix, meaning that ki j = k ji , where ki j are the components of k . We
will not go into the reason why just yet, but it is a quite general property of the stiffness
matrix (which you can use as a check on your algebra: if the stiffness matrix you have
found is not symmetric, then you have made a mistake!). The dynamical matrix D is also
symmetric, trivially so here as it is the same as k aside from the scalar multiplicative factor
of 1/m. We will make use of this symmetry property shortly.
To solve Eq. (4.78), we follow the same procedure we used for solving problems with
single oscillators; we look for solutions of the form
x = a e−iωt , (4.80)
where
a1
!
a= . (4.81)
a2
Substituting Eq. (4.80) into Eq. (4.78) yields
−ω2 a = −D a , (4.82)
which, upon rearranging can be written in the form of an eigenvalue equation
. (4.83)
92 Normal Modes
t
where λ = ω2 is called the eigenvalue and a is known as the eigenvector. Our task is to
determine the eigenvalues λ and associated eigenvectors a for the matrix D , which contain
all the essential physics of the problem. The eigenvectors will give us the normal frequen-
cies and the eigenvectors will give us the normal coordinates. In this case, there are two
eigenvalues and two eigenvectors, which follows from the fact that there are two degrees
of freedom, or equivalently that D is a 2 × 2 matrix.
Because D is a symmetric matrix, all its eigenvalues λ are guaranteed to be real. More-
√ it here, D is positive definite, meaning that its eigenvalues λ
over, while we do not prove
are positive. Since ω = λ, the normal frequencies ω are real. Physically, this is what we
expect, as complex frequencies would imply damping, which isn’t present in this problem.
Solving an eigenvalue problem

To solve the eigenvalue problem, we rewrite Eq. (4.83) as D a = λI a, where I is the identity
matrix defined by Ii j = δi j , and then subtract the right side from both sides of the equation:
(D − λI )a = 0 . (4.84)
Equation (4.84) has a solution only if the determinant of the matrix (D − λI ) is zero, that
is, if det(D − λI ) = 0 (see Appendix A). Writing out the determinant explicitly for this
problem gives
ω2a + ω2b − λ −ω2b

!
det(D − ω2I ) = det =0, (4.85)
−ω2b ω2b − λ
which gives
(ω2a + ω2b − λ)(ω2b − λ) − ω4b = 0 . (4.86)
Expanding Eq. (4.86) gives:
λ2 − (ω2a + 2ω2b ) λ + ω2a ω2b = 0 . (4.87)
The determinant of D is a second order polynomial called the characteristic polynomial

of the matrix D . We encountered this earlier in §4.1.1 (see Eq. (4.10)). In general, the
characteristic polynomial of an n × n matrix is an nth -order polynomial. Here it is a second
order polynomial, which we can solve using the quadratic formula.
Instead of solving Eq. (4.87) for the general √
case for arbitrary spring constants, we con-
sider the special case where kb /ka = ω2b /ω2a = 3/2, which simplifies the algebra, but still
serves to illustrate the method. With this choice, the solution to Eq. (4.87) is
√
 3 + 1 0.366 ω2a
 (
λ = ω2 =  ± 1 ω2a ≈ . (4.88)

2 2.366 ω2a
t
Thus, the two normal frequencies for this system are
√
s
3−1 √
ω1 = ωa ≈ 0.366 ωa ≈ 0.605 ωa (4.89)
2
√
s
3+3 √
ω2 = ωa ≈ 2.366 ωa ≈ 1.538 ωa . (4.90)
2
Associated with the normal frequencies ω1 and ω2 are their respective eigenvectors a1 and
√ the eigenvector a√1 associated with ω1 , we
a2 , which we can find using Eq. (4.84). To find
substitute a1 into Eq. (4.84) using λ1 = ω1 = 3−1
2 2 2 3 2
2 ωa and ωb = 2 ωa . This gives
ω2a + ω2b − λ −ω2b a1,1

! !
2

D − ω1I a1 = =0 (4.91)
−ω2b ω2b − λ a2,1
√
 3 ω2a − 23 ω2a 
 
a1,1
!
=  3 2
2√ =0, (4.92)
a2,1

1 2 
− 2 ωa 2 ωa

where the second subscript for a1,1 and a2,1 is a 1 to indicate that these are components of
the eigenvector for normal mode 1 (which has a normal frequency of ω1 . Multiplying out
the top and bottom rows yields
√
3 2 3 2
ω a1,1 − ω a2,1 = 0 (4.93)
2 a 2 a
√
3 2 1
− ω a1,1 + ω2a a2,1 = 0 . (4.94)
2 a 2
Solving either of these equations for a2,1 yields the same result:
√
a2,1 = 3 a1,1 . (4.95)
√
3+3
Repeating the same procedure but substituting a2 into Eq. (4.84) using λ2 = ω22 = 2 ω2a
yields two similar equations. Keeping only the equation from the top row gives
√
1 2 3 2
− ωa a1,2 − ω a2,2 = 0 . (4.96)
2 2 a
which yields
1
a2,2 = − √ a1,2 . (4.97)
3
If a1 and a2 , are eigenvectors, then any constant times a1 and a2 is also an eigenvector
(which, fundamentally, is why the above procedure determines only the ratios a2,1 /a1,1 and
a2,2 /a1,2 ). We will denote the normalized (length 1) eigenvectors as ẽ1 and ẽ2 , which are
given by
 1   √3 
 2  − 
ẽ1 =  √  , ẽ2 =  2  . (4.98)
3 1
2 2
94 Normal Modes
t
4.2.2 The meaning of normal coordinates
What is the physical interpretation of these normal coordinates we have found? To answer
this question, we return to Eq. (4.80), x = a e−iωt , and write it out explicitly for the first of
the two solutions we have found
x1,1 a
! !
= 1,1 e−iω1 t = ẽ1 A1 e−iω1 t (4.99a)
x2,1 a2,1
 1 
ẽ1,1 
 
 2 
=   A1 e 1 =  √  A1 e−iω1 t ,
−iω t
(4.99b)
ẽ2,1 3
2
where the first subscript on xi,µ and ai,µ refers to the displacement and amplitude of mass 1
or mass 2, and the second subscript (after the comma) designates the normal mode (here,
normal mode 1). These equations tell us that in normal √ mode 1, mass 1 and mass 2 al-
ways move with relative amplitudes of a2,1 /a1,1 = 3. The absolute overall scale of the
amplitude is determined by the initial conditions and set by the (complex) number A1 in
the above equations. The fact that both components of the eigenvector ẽ1 , namely ẽ1,1 and
ẽ2,1 , are positive means that in normal mode 1, both particles are always displaced from
equilibrium in the same direction. Moreover, in mode 1, both particles oscillate with the
same frequency, ω1 ≈ 0.605 ωa .
We can write similar equations for normal mode 2:
x1,2 a
! !
= 1,2 e−iω2 t = ẽ2 A2 e−iω2 t (4.100a)
x2,2 a2,2
 √3 
ẽ1,2 
 
− 
=   A2 e 2 =  2  A2 e−iω2 t ,
−iω t
(4.100b)
ẽ2,2 1
2
√
In mode two a2,2 /a1,2 = −1/ 3, meaning that the two masses are always displaced in
opposite directions, albeit with different amplitudes. In mode 2, both particles oscillate
with the same frequency, ω2 ≈ 1.538 ωa .
You should think of the normal coordinates as giving the relative amplitudes of all the
particles in a given normal mode. All particles in that mode oscillate with the same fre-
quency, the normal frequency, also known as the eigenfrequency.
The complete solution to the normal mode problem is expressed as the sum of the two
normal modes
x = A1 ẽ1 e−iω1 t + A2 ẽ2 e−iω2 t (4.101)

N
X
= An ẽn e−iωn t , (4.102)
n=1
where, for this case, the number of normal modes N = 2. The amplitudes and phases of A1
and A2 are determined by the initial conditions.
t
t
Fig. 4.6 Normal coordinates for 2-mass 2-spring problem. Unit vectors e1 and e2 are shown
along the x1 and x2 directions (gray) and for the normal coordinates ẽ1 and ẽ2 (black).
Transforming equations to normal coordinates

Figure 4.6 shows the eigenvectors ẽ1 and ẽ2 graphically, plotted in a basis (coordinate
system) defined by the particle coordinates (x1 , x2 ). Note that ẽ1 and ẽ2 are perpendicular
to each other and can be written as
cos θ − sin θ

   
ẽ1 =    , ẽ2 = 
   , (4.103)
sin θ cos θ
where θ = 60◦ for this problem. Other choices of the ratio kb /ka give similar results, with ẽ1
and ẽ2 perpendicular to each other and rotated with respect to e1 and e2 , but with different
angles θ. Note that we have chosen the signs for the components e1 and e2 to give a right
handed coordinate system. In particular, this meant choosing the upper component of the
e2 vector to be negative. Making such a choice is not strictly necessary, but it does make it
somewhat simpler to transform results between the different coordinate systems.
The vectors ẽ1 and ẽ1 define a new coordinate system that is rotated by the angle θ with
respect to the original coordinate system defined by the unit vectors e1 and e2 along the x1
and x2 axes, respectively:
1 0
! !
e1 = , e2 = . (4.104)
0 1
The solutions to the problem can be expressed in either set of coordinates: in the {e1 , e2 }
coordinate system, solutions are expressed as functions of the original coordinates x1 and
x2 ; in the {ẽ1 , ẽ2 } coordinate system, solutions are expressed as functions of the rotated
coordinates x̃1 and x̃2 . In the context of linear algebra, each of these coordinate systems
is referred to as a basis. We have constructed each basis to be orthonormal, meaning that
all the basis vectors are normalized to have a length of 1, and each set of basis vectors are
orthogonal to each other: e1 · e2 = 0 and ẽ1 · ẽ2 = 0.
Any vector x describing the positions of the masses can be expressed either basis:
x = x1 e1 + x2 e2 (4.105a)
= x̃1 ẽ1 + x̃2 ẽ2 . (4.105b)
96 Normal Modes
t
For example, we can express the components x̃1 and x̃2 in either basis using dot products:
x̃1 = x · ẽ1 = (x1 e1 + x2 e2 ) · ẽ1 = (e1 · ẽ1 ) x1 + (e2 · ẽ1 ) x2 (4.106a)
x̃2 = x · ẽ2 = (x1 e1 + x2 e2 ) · ẽ2 = (e1 · ẽ2 ) x1 + (e2 · ẽ2 ) x2 . (4.106b)
We can rewrite Eqs. (4.106a) and (4.106b) as a single matrix equation
x̃ = Ux , (4.107)
where
x̃1 x1
! !
x̃ = , x= , (4.108)
x̃2 x2
and
cos θ sin θ
! !
e1 · ẽ1 e2 · ẽ1
U= = (ẽ1 , ẽ2 ) = . (4.109)
e1 · ẽ2 e2 · ẽ2 − sin θ cos θ
Note that the columns of U are the eigenvectors ẽ1 and ẽ2 we found in Eqs. (4.98) and
(4.103). The matrix U is a rotation matrix that transforms (rotates) coordinates x in the
original (e1 , e2 ) basis to coordinates x̃ in the eigenvector (ẽ1 , ẽ2 ) basis.
The equation of motion in matrix form is given by
ẍ = −D x . (4.110)
We can rewrite Eq. (4.110) in the tilde basis. Noting from Eq. (4.107) that x = U −1 x̃, Eq.
(4.110) can be written as
U −1 x̃¨ = −DU −1 x̃ . (4.111)
Left multiplying both sides by U , this becomes
x̃¨ = −UDU−1 x̃ , (4.112)
or
x̃¨ = −D̃
Dx̃ , (4.113)
where
D ≡ UDU −1 = UDU T ,
D̃ (4.114)
where we have noted the property of rotation matrices that their inverse is their transpose.
The two eigenvectors, with eigenvalues λ1 = ω21 and λ2 = ω22 , point along the ẽ1 and ẽ2
axes, respectively, so within the tilde basis,
1 0
! !
ẽ1 = , ẽ2 = . (4.115)
0 1
D ã = λ ã, D̃
Since D̃ D must be a diagonal matrix with the eigenvalues as the diagonal elements
λ1 0 ω21 0
! !
D̃ ≡ UDU T = = . (4.116)
0 λ2 0 ω22
t
Writing Eq. (4.113) in component (scalar) form with x̃1 = x̃1 ẽ1 and x̃2 = x̃2 ẽ2 , we thus
obtain two uncoupled equations of motion
x̃¨ µ = −D̃
Dx̃µ
x̃¨µ ẽµ = −D̃
D x̃µ ẽµ = −λµ x̃µ ẽµ ⇒ x̃¨µ = −ω2µ x̃µ , (4.117)
with µ = 1, 2. Thus, our procedure has reduced the original two coupled equations of
motion to two uncoupled equations of motion.
4.2.3 Specific solutions
To find solutions for a specific set of initial conditions, we proceed the same way we pro-
ceeded in Example 4.1, remembering that we are now using the notation x̃ in place of q to
refer to the normal coordinates. Thus, as before, since x̃1 and x̃2 obey equations of motion
for uncoupled simple harmonic oscillators, the solutions can be written in any one of the
usual forms. Here we use sines and cosines
x̃1 (t) = A1 cos ω1 t + B2 sin ω1 t (4.118)

x̃2 (t) = A2 cos ω2 t + B2 sin ω2 t , (4.119)
and determine the four amplitudes from the initial conditions, which are generally ex-
pressed in terms of the original variables. For example, suppose x1 (0) = 4 and x2 (0) = 2
and ẋ1 (0) = ẋ2 (0) = 0, which means that
4 0
! !
x= , ẋ = . (4.120)
2 0
We can use Eq. (4.107) to find the initial conditions for the normal coordinates
x̃1 (0) cos θ sin θ x1 (0)

! ! !
x̃(0) = = Ux = (4.121)
x̃2 (0) − sin θ cos θ x2 (0)
√  
3
√  
 1

 4  2 + 3   3.732 
 
2 2 
=  √    =  √  '   , (4.122)
1  2
− 23 2
1 − 2 3 −2.464
where θ = 60◦ . For the initial velocities, we have
˙
˙x̃(0) = x̃1 (0) = U ẋ = cos θ sin θ ẋ1 (0)
! ! !
(4.123)
x̃˙2 (0) − sin θ cos θ ẋ2 (0)
√  
 1 3

 0 0
 
2 2 
=  √    =   , (4.124)
1  0 0
− 23 2
Applying these initial conditions to Eqs. (4.118) and (4.119) gives A1 ' 3.732, B1 = 0,
A2 ' −2.464, and B2 = 0. To find solutions for x1 (t) and x2 (t), we use the inverse of Eq.
98 Normal Modes
t
(4.107)
x(t) = U −1 x̃(t) = U T x̃(t) (4.125)

√
3 
 x1 (t)  12
  
  x̃1 (t)
 
− 2 
  =  √   , (4.126)
x2 (t) x̃2 (t)
 3  
1
2 2
where we have used the property of rotation matrices that U −1 = U T . With these results,
we obtain
√
3
x1 (t) = 21 x̃1 (t) − 2 x̃2 (t) (4.127)
√
1 3
' 2 3.732 cos ω1 t − 2 (−2.464 cos ω2 t) (4.128)
' 1.866 cos ω1 t + 2.1344 cos ω2 t . (4.129)
and
√
3 1
x2 (t) = 2 x̃1 (t) + 2 x̃2 (t) (4.130)
√
3
' 2 3.732 cos ω1 t + 21 (−2.464 cos ω2 t) (4.131)
' 3.232 cos ω1 t − 1.232 cos ω2 t . (4.132)
Thus, with Eqs. (4.129) and (4.132) our solution is complete.
4.2.4 Three masses on a string
We would like to find the normal modes for a system consisting of three equal masses
attached to a massless string stretched between two fixed walls, as depicted in Fig. 4.7.
The string is held under tension T with the masses equally spaced at distance l from each
other and from the walls. Suppose one or more of the masses is pulled up and then released
setting the system into oscillation in the vertical plane. The vertical distance of each mass
from its equilibrium position is yi , where i = 1, 2, 3. We take the amplitude of the oscilla-
tions to be small compared to the distance between masses so that yi l at all times.
The first step is to write down the equations of motion for each mass. Let’s begin with
the leftmost mass whose vertical position is y1 . The vertical force on the mass will be the
difference in the vertical component of the tension of the string on either side of the mass
F1y = −T sin θ0 + T sin θ1 (4.133)
t
Fig. 4.7 Three masses on a string under tension. The vertical scale is greatly magnified
compared to the horizontal.
t
In the limit of small amplitude oscillations where yi l, the sines can be approximated as
sin θ0 ' y1 /l and sin θ1 ' (y2 − y1 )/l. Thus, the force on the leftmost mass is
y1 y2 − y1 T
F1y = −T +T = (−2y1 + y2 ) (4.134)
l l l
Using similar analysis, the vertical forces on the middle and rightmost masses are
y2 − y1 y3 − y2 T
F2y = −T +T = (y1 − 2y2 + y3 ) (4.135)
l l l
y2 − y3 y3 T
F3y =T −T = (y2 − 2y3 ) (4.136)
l l l
You should check that the signs of the vertical forces are in the correct directions for the
configurations of the masses used in Fig. 4.7. We can now write down the equations of
motion of each of the three masses:
T
mÿ1 = (−2y1 + y2 ) (4.137)
l
T
mÿ2 = (y1 − 2y2 + y3 ) (4.138)
l
T
mÿ3 = (y2 − 2y3 ) . (4.139)
l
The dynamical matrix for this system of equations is thus given by
 2 −1 0 
 
D = ω20 −1 2 −1 , (4.140)

 
0 −1 2
 
√
where ω0 = T /ml. The characteristic frequency ω0 that emerges increase when the
tension in the string is increased or when the mass or distanced between the masses is
increased, which makes sense physically.
The normal frequencies are found by finding the solutions to the characteristic equation,
which is given by det(D − λI ) = 0 where λ = ω2 . Defining s = λ/ω20 , this becomes
2 − s −1 0 
 
det(D − λI ) = ω20 det  −1 2 − s −1  = 0 (4.141)

 
0 −1 2 − s
 
2 3
h i
= ω0 (2 − s) − 2(2 − s) = 0 , (4.142)
√
which has solutions of s = 2 and s = 2 ± 2. Therefore, the normal frequencies, from the
lowest to the highest frequency, are given by
√
q
ω1 = 2 − 2 ω0 ≈ 0.765 ω0 (4.143)
√
ω2 = 2 ω0 ≈ 1.414 ω0 (4.144)
√
q
ω3 = 2 + 2 ω0 ≈ 1.848 ω0 , (4.145)
√
where ω0 = T /ml.
100 Normal Modes
t
The eigenvectors are found by solving (D − ω2µI )a = 0 for each normal frequency ωµ
with µ = 1, 2, 3. Starting with ω1 , we have (D − ω21I )a1 = 0, which written out gives
√
2 − (2 − 2) 0  a1,1 
  
−1
√
 2 (2 2) a =0 (4.146)
  
−1 − − −1
 √   2,1 
0 −1 2 − (2 − 2) 3,1 a
Multiplying out the three rows, we obtain

√
2 a1,1 − a2,1 = 0 (4.147)
√
−a1,1 + 2 a2,1 − a3,1 = 0 (4.148)
√
−a2,1 + 2 a3,1 = 0 . (4.149)
√ √
From the first and third equations we have a2,1 = 2 a1,1 and a2,1 = 2 a3,1 , which means
that a3,1 = a1,1 . The second equation is automatically satisfied by these relations. Thus, the
first normalized eigenvector is
 1 
 
1  √ 
ẽ1 =  2 . (4.150)
2  
1
Next, we solve (D − ω22I)a2 = 0, which gives
2 − 2 −1 0  a1,2 

  
 −1 2 − 2 −1  a2,2  = 0 (4.151)

  
0 −1 2 − 2 a3,2
  
Again, multiplying out the three rows, we obtain
−a2,2 = 0 (4.152)
−a1,2 − a3,2 = 0 (4.153)
−a2,2 = 0 . (4.154)
The first and third equations give a2,2 = 0. The second equation gives a1,2 = −a3,2 . Thus,
the second normalized eigenvector is
1
 
1  
ẽ2 = √  0  . (4.155)
2 −1
Finally, we solve (D − ω23I )a3 = 0, which gives

√
2 − (2 + 2) 0  a1,3 
  
−1
√
 2 − (2 + 2) a =0 (4.156)
  
−1 −1
 √   2,3 
0 −1 2 − (2 + 2) a3,3
t
Multiplying out the three rows, we obtain
√
− 2 a1,1 − a2,1 = 0 (4.157)
√
−a1,1 − 2 a2,1 − a3,1 = 0 (4.158)
√
−a2,1 − 2 a3,1 = 0 . (4.159)
√ √
From the first and third equations we have a2,1 = − 2 a1,1 and a2,1 = − 2 a3,1 , which
means that a3,1 = a1,1 . The second equation is automatically satisfied by these relations.
Thus, the third normalized eigenvector is
 1 
 
1 √
ẽ3 = − 2 . (4.160)

2 
1

t
Fig. 4.8 Normal modes of 3 masses on a string. The relative amplitudes in each normal mode
are proportional to the components of the corresponding eigenvector. The
eigenvectors and normal frequencies are shown to the left and right, respectively, of
each drawing.
Figure 4.8 shows the normal modes for three masses on a string under tension. The three
components of each of the three eigenvectors give the relative amplitudes of the masses in
each normal mode.
Example 4.2 Consider the case of three masses on a string where m = 25 g, l = 10 cm,
and T = 50 N. Suppose that the middle mass is given an initial velocity of 30 = 0.50 m/s,
while the other two masses are initially stationary. All three masses start out in their equi-
librium positions. What is the subsequent motion of the masses? What frequencies are
produced?
102 Normal Modes
t
Solution
Let’s start by calculating the characteristic vibration frequency ω0 :
r s
T 50 N
ω0 = = = 447 s−1 (4.161)
ml (0.025 kg)(0.010 m)
ω0
f0 = = 71.2 Hz . (4.162)
2π
The initial conditions are
y1 = y2 = y3 = 0 , 31 = 33 = 0 , 33 = 30 = 0.50 m/s (4.163)
Note that with these initial conditions, masses 1 and 3 must have identical trajectories,
because they are symmetrically placed about mass 2, which is the only mass that is initially
perturbed from equilibrium. This means that normal mode 2 will not be excited in this case,
because in normal mode 2, masses 1 and 3 are displaced in opposite directions. Another
way of looking at this is to note that the initial conditions are symmetric with respect to the
central particle, as are normal modes 1 and 3. By contrast, normal mode 2 is antisymmetric
with respect to mass 2, so it will not be excited in this case.
Generalizing Eqs. (4.106a) and (4.106b) to a 3-particle system, the normal coordinates
are
1h √ i
ỹ1 (t) = y · ẽ1 = y1 (t) + 2 y2 (t) + y3 (t) (4.164)
2
1
ỹ2 (t) = y · ẽ2 = y1 (t) − y3 (t) (4.165)

2
1 h √ i
ỹ3 (t) = y · ẽ3 = y1 (t) − 2 y2 (t) + y3 (t) . (4.166)
2
Let’s determine the initial conditions for the normal coordinates.
ỹ1 = y˜2 = ỹ3 = 0 (4.167)

1h √ 1
ỹ˙ 1 =
i
ẏ1 (0) + 2 ẏ2 (0) + ẏ3 (0) = √ 30 (4.168)
2 2
1
ỹ˙ 2 = √ ẏ1 (0) − ẏ3 (0) = 0 (4.169)

2
1 √ 1
ỹ˙ 3 =
h i
ẏ1 (0) − 2 ẏ2 (0) + ẏ3 (0) = − √ 30 . (4.170)
2 2
We will write the solutions to each normal mode as a sum of a sine and a cosine function.
The general solutions for the different normal modes are
ỹ1 (t) = A1 cos ω1 t + B1 sin ω1 t (4.171)

ỹ2 (t) = A2 cos ω2 t + B2 sin ω2 t (4.172)
ỹ3 (t) = A3 cos ω3 t + B3 sin ω3 t . (4.173)
Because ỹ1 = ỹ2 = ỹ3 = 0, A1 = A2 = A3 = 0. Similarly, because ỹ˙ 2 = 0, B2 = 0, thus

verifying our earlier claim that normal mode 2 would not appear in the solution for this set
t
of initial conditions. Applying the remaining initial conditions
1
ỹ˙ 1 = ω1 B1 = √ 30 (4.174)
2
1 30 1 30
⇒ B1 = √ = q √ = 1.03 mm (4.175)
2 ω1 2 2− 2 0
ω

1
ỹ˙ 3 = ω3 B3 = √ 30 (4.176)
2
1 30 1 30
⇒ B3 = √ = q √ ω = 2.92 mm . (4.177)
ω
2 3 2 2+ 2 0
Note that B1 l and B3 l, justifying our approximation that the deflections of the string
are small compared to the spacing between particles. Because ỹ2 (t) = 0, y1 (t) = y3 (t).
Adding and subtracting Eqs. (4.164) and (4.166) thus gives
1 1
y1 (t) = y3 (t) = ỹ1 (t) + ỹ3 (t) = [B1 sin ω1 t + B3 sin ω3 t] (4.178)
2 2
1 1
y2 (t) = √ ỹ1 (t) − ỹ3 (t) = √ [B1 sin ω1 t − B3 sin ω3 t] . (4.179)

2 2
4.2.5 Systems of unequal masses
The analytical method we have developed thus far is suited only for coupled systems of
equal masses. When the masses are not equal, forming the dynamical matrix becomes a bit
more involved.
Let’s consider the problem of two pendulums coupled by a spring with a force constant k
that we introduced at the beginning of this chapter, but this time let’s make the two masses,
denoted M and m, different. Figure 4.1 still serves to define the problem, and aside from
the different masses, the equations of motion are the same as they were before
Mg
M ẍ1 = − + k x1 + kx2 (4.180)
l
mg
m ẍ2 = kx1 − + k x2 , (4.181)
l
where x1 and x2 are the coordinates of the masses M and m, respectively. We can write
these equations in matrix form
m ẍ = −k x , (4.182)
where
M 0 Mg/l + k x1
! ! !
−k
m= , k= , x= . (4.183)
0 m −k mg/l + k x2
We wish to recast Eq. (4.182) in the standard form ẍ = −D x. The most straightforward
104 Normal Modes
t
approach is simply to left multiply both sides of Eq. (4.182) by the inverse of the mass
matrix m −1 , which gives
ẍ = −D 0 x (4.184)
where
M −1 0 Mg/l + k
! !
−k
D =m k =
0 −1
(4.185)
0 m−1 −k mg/l + k
g/l + k/M ω2 + ω2 −ω2s
! !
−k/M
= = p 2s , (4.186)
−k/m g/l + k/m −ηω s ω2p + ηω2s
√
where ω p = g/l, ω s = k/M, and η = M/m. The eigenvalues λ of D 0 are given by
p
solving det(D 0 − λI ) = 0:
ω2p + ω2s − λ −ω2s
!
det = (ω2p + ω2s − λ)(ω2p + ηω2s − λ) − ηω4s = 0 (4.187)
−ηω2s ω2p + ηω2s − λ
= (λ − ω2p )(λ − ω2p − (1 + η)ω2s ); . (4.188)
Solving for λ = ω2 gives
2

ω p

λ = ω2 =  , (4.189)

ω2 + (1 + η)ω2

p s
which reduces to the result we obtained in §4.1.1 for equal masses when η = 1.
We could go on to find the eigenvectors, which proceeds along now familiar lines. How-
ever, the dynamical matrix D 0 is not symmetric, and this has certain disadvantages. First,
when the dynamical matrix is not symmetric, the eigenvectors do not in general form an
orthogonal system. As we shall see, the can be a disadvantage in performing certain calcu-
lations. Second, for systems with a large number of degrees of freedom, the eigenvalues and
eigenvectors must be found numerically in most cases. Because a non-symmetric matrix
has nearly twice as many distinct entries as its symmetric counterpart for which Di j = D ji ,
numerical computation of eigenvalues and eigenvectors require more memory and take
considerably more time for a non-symmetric matrix than for its symmetric counterpart.
For the diagonal mass matrices m we have encountered here, it is a simple matter to
rewrite the eigenvalue equation so that the dynamical matrix is diagonal. We start by defin-
ing the square root m 1/2 and inverse square root m −1/2 of the mass matrix
M 1/2 0 M −1/2 0
! !
1/2
m = , m −1/2
= . (4.190)
0 m1/2 0 m−1/2
Starting from m ẍ = −k x, Eq. (4.182), we left multiply both sides by m −1/2

m −1/2m ẍ = −m −1/2k m −1/2m 1/2 x ,

(4.191)
and insert the identity matrix I = m−1/2m1/2 as shown on the right hand side. Regrouping
terms, this becomes
m 1/2 ẍ = − m −1/2k m −1/2 m 1/2 x ,

(4.192)
t
or
ÿ = −D y , (4.193)
where
√ √
g/l + k/M −k/ Mm ω2 + ω2 − η ω2s
! !
D =m −1/2
km −1/2
= √ = p√ 2s , (4.194)
−k/ Mm g/l + k/m − η ωs ω2p + ηω2s
is now a symmetric dynamical matrix and

√
M 1/2 0 x1 M x1
! ! !
1/2
y=m x= = √ . (4.195)
0 m1/2 x2 m x2
The only price we have to pay for this transformation is the introduction of the vector y. A
quick examination of D 0 and D reveals that they have the same characteristic polynomial,
and thus the same set of eigenvalues. However, they do not have the same eigenvectors.
To find the eigenvectors, we need to solve (D − ω2µI )a = 0 for each frequency ωµ . While
we can perform the calculation for arbitrary ω p and ω s , the algebra is messy. To simplify
the algebra, we choose ω p = 6, ω s = 2, and η = M/m = 9/4. For this case, Eq. (4.194)
gives
40 −6
!
D= (4.196)
−6 45
while Eq. (4.189) gives λ = 36 and 49 or ω1 = 6 and ω2 = 7. The eigenvectors are given
by solving (D − ω2µI )a = 0 for ω1 = 6 and ω2 = 7:
40 − 36 a1,1 40 − 49 a1,2
! ! ! !
−6 −6
=0, =0, (4.197)
−6 45 − 36 a2,1 −6 45 − 49 a2,2
which gives normalized eigenvectors of
1 3 cos θ 1 − sin θ
! ! ! !
−2
ẽ1 = √ = , ẽ2 = √ = , (4.198)
13 2 sin θ 13 3 cos θ
where as before the the normalized two-component eigenvectors can be written in terms of
sines and cosines of the rotation angle θ; in this case θ ' 33.7◦ .
We can now use these results to solve an initial value problem. Let’s take initial condi-
tions of x1 (0) = x0 , x2 (0) = 0, and ẋ1 (0) = ẋ2 (0) = 0. We need to express these in terms of
y(0) and ẏ(0), which we do using Eq. (4.195)
√ √
y1 (0) M x (0) M x0
! ! !
y(0) = = √ 1 = (4.199)
y2 (0) m x2 (0) 0
√
ẏ1 (0) M ẋ1 (0) 0
! ! !
ẏ(0) = = √ = . (4.200)
ẏ2 (0) m ẋ2 (0) 0
The general solutions, written in terms of the normal coordinates ỹ1 (t) and ỹ2 (t), are given
106 Normal Modes
t
by the usual sine and cosine functions
1
ỹ1 (t) = y · ẽ1 = √ [3y1 (t) + 2y2 (t)] = A1 cos ω1 t + B1 sin ω1 t (4.201)
13
1
ỹ2 (t) = y · ẽ2 = √ [−2y1 (t) + 3y2 (t)] = A2 cos ω2 t + B2 sin ω2 t . (4.202)
13
The initial conditions on ỹ1 and ỹ2 are
r
1 M
ỹ1 (0) = √ 3y1 (0) + 2y2 (0) = 3x0 = A1 (4.203)

13 13
r
1 M
ỹ2 (0) = √ −2y1 (0) + 3y2 (0) = − 2x0 = A2 (4.204)

13 13
ỹ˙ 1 (0) = ỹ˙ 2 (0) = 0 . (4.205)
The initial conditions on the velocities ỹ˙ 1 (0) and ỹ˙ 2 (0) give B1 = B2 = 0. Thus, the solu-
tions are given by
r r
M M
ỹ1 (t) = 3x0 cos ω1 t , ỹ2 (t) = − 2x0 cos ω2 t . (4.206)
13 13
√
Solving Eqs. (4.201)
√ and (4.202) for y1 and y2 , we obtain y1 (t) = (1/ 13)[3ỹ1 (t) − 2ỹ2 (t)]
and y2 (t) = (1/ 13)[2ỹ1 (t) + 3ỹ2 (t)]. Multiplying Eq. (4.195) by m −1/2 gives
 1   1
 √ y1   √ [3ỹ1 (t) − 2ỹ2 (t)]

M 13M
x=m −1/2
y =  1  =  1  . (4.207)
√ y2
m
√
13m
[2ỹ 1 (t) + 3ỹ2 (t)]
Substituting the expressions for ỹ1 (t) and ỹ2 (t) from Eq. (4.206) into Eq. (4.207) gives
9 4
!
x1 (t) = x0 cos ω1 t + cos ω2 t (4.208)
13 13
√
6 η
x2 (t) = x0 (cos ω1 t − cos ω2 t) . (4.209)
13
A quick check on these solutions reveals that they satisfy the initial conditions that x1 (0) =
x0 , x2 (0) = 0, and ẋ1 (0) = ẋ2 (0) = 0. They also satisfy the equations of motion, as must
any solution that is a sum sine and cosine functions of ω1 t and ω2 t.
When ω1 and ω2 are nearly equal, as they are in this case where ω1 = 6 and ω2 = 7,
we expect beats to appear. In this case, it is useful to rewrite the solutions in terms of the
average and difference frequencies, ω̄ ≡ 21 (ω1 + ω2 ) and ∆ω ≡ 21 (ω2 − ω1 ), respectively.
The normal mode frequencies are then written as ω1 = ω̄ − ∆ω and ω2 = ω̄ + ∆ω. Using
cosine addition identities, the solutions can be rewritten as
5
!
x1 (t) = x0 cos ∆ωt cos ω̄t + sin ∆ωt sin ω̄t (4.210)
13
√
12 η
x2 (t) = x0 sin ∆ωt sin ω̄t (4.211)
13
The trajectories for the two pendulums are shown in Fig. 4.9. As in the case treated
previously for equal-mass coupled pendulums, a beating pattern is observed. In contrast
t
x1 (t)
1
0 t
2 4 6 8 10 12 14
−1
x2 (t)
0 t
2 4 6 8 10 12 14
−1
t
Fig. 4.9 Incomplete modulation of trajectories of weakly-coupled pendulums of unequal
masses. The dashed lines show that time-dependent amplitudes of the oscillations.
to the previous case, however, the amplitude of the oscillations of the heavier mass goes
through a minimum but never goes to zero. More generally, one can show3 that
x1 = x0 cos ∆ωt cos ω̄t + (cos2 θ − sin2 θ) sin ∆ωt sin ω̄t ,

(4.212)
which means that the case where x1 (t) goes to zero when x2 (t) is a maximum is a special
case that only occurs when cos2 θ − sin2 θ = 0 as is does, for example, when θ = 45◦ . This
is precisely what occurs for the case when M = m, which we studied in §4.1.1.
4.2.6 Geometry and symmetry
The transformation to a coordinate system in which D is symmetric may seem like a lot
more trouble than it is worth. In some cases that may be true, but there are practical ben-
efits that may not be readily apparent. The most important one is that the eigenvectors are
orthogonal when D is symmetric, and many other useful properties follow from this one.
Let’s take a closer look at the problem of the pendulum with unequal masses that we
just solved. Once again we consider the particular case in which ω p = 6, ω s = 2, and
η = M/m = 9/4, and calculate the eigenvectors of D 0 in the original basis. In that basis, the
3
√
√ the more general result, redo the calculation of starting with Eq. (4.201) and let (2/ 13) → sin θ
To obtain
and (3/ 13) → cos θ.
108 Normal Modes
t
t
Fig. 4.10 Coordinates for a coupled pendulum with unequal masses. (a) Normal coordinates ẽ01
in (x1 , x2 ) coordinates: α ' 69◦ and θ0 = 45◦ . (b) Normal coordinates in stretched
(y1 , y2 ) coordinates: α = 90◦ and θ ' 33.7◦ .
dynamical matrix is given by

40
!
−4
D =
0
, (4.213)
−9 45
which is not symmetric. The eigenvalues of D 0 are λ = 36 and 49 or ω1 = 6 and ω2 = 7,
the same as for D , as they must be. But the eigenvectors are different. With a little algebra,
the normalized eigenvectors of D0 are found to be
1 1 1 −4
! !
ẽ1 = √
0
, ẽ2 = √
0
. (4.214)
2 1 97 9
The cosine of the angle α0 between e01 and e02 is given by their dot product, from which we
find that α0 ' 69◦ : the eigenvectors are not orthogonal. Figure 4.10(a) shows the geometry
of the eigenvectors ẽ01 and ẽ01 .
Figure 4.10(b) illustrates how the transformation to the y basis, given by Eq. (4.195),
stretches the (x1 , x2 ) coordinate system so that the basis vectors become orthogonal, with
α = 90◦ in the new basis. After stretching the basis to obtain orthogonal normal coordi-
nates, the coordinates are renormalized to have unit length once again.
Thus, while the ostensible purpose of the transformation to the y basis was to obtain a
symmetric dynamical matrix, from the standpoint of geometry, the purpose of the trans-
formation is to find a coordinate system in which the eigenvectors are orthogonal. It is this
orthogonality that makes it possible to obtain the normal coordinates as a pure rotation and
to express the eigenvectors in terms of sines and cosines of the rotation angle. This, in turn,
allows us to obtain the result given by Eq. (4.212), which helps us understand how beats
can appear either fully modulated, as seen in Fig. 4.4, or only partially modulated, as seen
in Fig. 4.9.
Formulating normal modes in the language of linear algebra and matrices brings a little
more order to the study of the subject. It also allows us to draw on a vast trove of math-
ematical theorems from linear algebra, which simplifies the analysis of such problems, as
well as giving physical insight. Thus far, we have only studied systems with a few degrees
of freedom. In Chapter 5 we take up the problem of normal modes again and apply matrix
methods to systems with a large number of degrees of freedom. For systems with many de-
grees of freedom, transforming the coordinates to make the dynamical matrix symmetric is
109 Normal modes of carbon dioxide
t
t
Fig. 4.11 Geometry for vibrations of triatomic CO2 molecule: (a) longitudinal, (b) transverse.
particularly important, as it is generally much more efficient computationally for systems

with many degrees of freedom.
4.3 Normal modes of carbon dioxide
Carbon dioxide is a linear molecule consisting of a single carbon atom sandwiched be-
tween two oxygen atoms. If we consider only small oscillations about the equilibrium
configurations of the three atoms, we can treat the problem of molecular vibrations within
the framework of normal modes that we have developed in this chapter.
4.3.1 Longitudinal modes
To begin, we consider only longitudinal motion, that is, motion along the axis formed by
the three atoms. Later we will consider tranverse motion, which is motion perpendicular to
the molecular axis. The longitudinal deviations of the three particles from their equilibrium
positions are denoted x1 , x2 , and x3 , as shown in Fig. 4.11(a). The mass of the two oxygen
atoms is m and the mass of the carbon atom is M. The equations of motion for the three
particles are
F1 = k(x2 − x1 ) = m ẍ1 (4.215)
F2 = k(x3 − x2 ) − k(x2 − x1 ) = M ẍ2 (4.216)
F3 = k(x2 − x3 ) = m ẍ3 , (4.217)
where the spring constant is given by
d2 U

k= , (4.218)
dξ2 ξ=0

U(ξ) is the interatomic potential between the oxygens and the carbon atom and ξ = x2 − x1
or ξ = x3 − x2 . Following the usual procedure and writing xi (t) = ai e−iωt for i = 1, 2, 3, we
obtain Eq. (4.78), which we repeat here for easy reference,
m ẍ = −k x , (4.219)
with the mass and stiffness matrices obtained from Eqs. (4.215)-(4.217) given by
110 Normal Modes
t
m 0 0   k 0 
   
−k
m =  0 M 0  , k = −k 2k −k . (4.220)
   
0 0 m 0 k
 
−k
 
The inverse of m is
m 0 0 
 −1 
m −1 =  0 M −1 0  , (4.221)

 
0 0 m−1
 
so that the dynamical matrix is given by

 k/m 0 
 
−k/m
D = m k = −k/M
−1
2k/M −k/M  , (4.222)
 
0 k/m

−k/m

which although not symmetric, will be used as is. Once again, the solution to the systems of
equations is given by solving the secular equation det(D − λI ) = 0, which for this problem
is
k/m − λ 0 
 
−k/m
det(D − λI ) = det  −k/M 2k/M − λ −k/M  = 0 , (4.223)

0 −k/m k/m − λ
 
where λ = ω2 . Evaluating the determinant, we obtain

k 2k k 2k2
!" ! ! #
−λ −λ −λ − =0, (4.224)
m M m mM
which upon simplification yields
k k 2m
!" !#
λ −λ λ− 1+ =0. (4.225)
m m M
Therefore, the normal frequencies are
r s
k k 2m
!
ωα = 0 ωβ = ωγ = 1+ , (4.226)
m m M
which corresponds to three normal modes and six solutions ±ωα , ±ωβ , and ±ωγ . At first
glance, the slowest mode with ωα = 0 might seem to be a mistake, until we recall that the
equation of motion for a normal mode is always of the form
q̈α + ω2α qα = 0 . (4.227)
In this case where ωα = 0, the equation of motion becomes
q̈α = 0 , (4.228)
which has the general solution qα (t) = qα (0) + 3α t, or motion along a straight line an
constant velocity 3α . Thus, one of our normal modes is not a vibrational mode at all but
corresponds to constant velocity rigid body motion of the center of mass.
Next we proceed to obtain the eigenvectors for each normal mode. This will reveal the
t
nature of the each of the modes corresponding to the normal frequencies ωα , ωβ , and ωγ .
Starting with the eigenvalue equation Eq. (4.83)
Da = λa ,
we set λ = ω2α = 0 to examine the α normal mode. For this case, Eq. (4.83) becomes
 k/m 0  a1  0

    
−k/m
2k/M −k/M  a2  = 0 . (4.229)
    
−k/M
0 −k/m k/m a3 0
    
Multiplying out the top row gives a1 − a2 = 0, which means that a1 = a2 . Multiplying
out the bottom row gives a3 − a2 = 0, which means that a3 = a2 . Thus, the normalized
eigenvector for the α mode is aTα = √13 (1, 1, 1), which is consistent with our earlier finding
that all three particles move together in a zero-frequency mode corresponding to rigid
translation.
To find the second normal mode, we set λ = ω2β = k/m. Equation (4.83) becomes
 k/m 0  a1  a 

    
−k/m
   k  1 
2k/M −k/M  a2  = a2  . (4.230)

−k/M
  m  
0 −k/m k/m a3 a3

Multiplying out the top row gives a1 −a2 = a1 , which means that a2 = 0. Setting a2 = 0 and
multiplying out the middle row gives a3 = −a1 . Thus, the normalized eigenvector for the
β mode is aTβ = √12 (1, 0, −1). This means that the β mode is one in which the two oxygen
atoms move exactly out of phase with each other and the central carbon atom remains
fixed, as does the center of mass of the molecule. Because the carbon atom does not move,
√
its mass M does not figure in the normal frequency ωβ = k/m for this problem.
To find the third normal mode, we set λ = ω2γ = (k/m)(1 + 2m/M). Equation (4.83)
becomes
 k/m 0  a1  " !# a1 

    
−k/m
k 2m  
2k/M −k/M  a2  = 1+ a2  . (4.231)
  
−k/M
m M  
0 −k/m k/m a3 a3
  
Multiplying out the top row gives a2 = −(2m/M)a1 . Similarly, multiplying out the bottom
row gives a2 = −(2m/M)a3 , which meansp that a1 = a3 . Thus the normalized eigenvector
for the γ mode is aTγ = (1, −2m/M, 1)/ 2 + 4(m/M)2 ).
To summarize, the three normal modes are characterized by the following eigenvectors
112 Normal Modes
t
and eigenvalues (or eigenfrequencies or normal frequencies)
1
 
1  
aα = √ 1 , λ ≡ ω2α = 0 (4.232)
3 1
 
−1
1   k
aβ = √  0  , λ ≡ ω2β = (4.233)
21   m
 1 
 
1 k 2m
!
aγ = p −2m/M  , λ ≡ ω2γ = 1+ (4.234)
2 + 4(m/M) )  1 
2
  m M
Figure 4.12(a) depicts the displacements of the atoms for each normal mode. Note that the
relative amplitudes of the displacements correspond to the relative amplitudes of the three
components of the corresponding eigenvector. This is a general feature of the eigenvectors
and thus provides a simple connection between the components of components of
The reason we obtained three normal modes from our analysis is because we allowed
for three degrees of freedom. We could have eliminated the trivial zero-frequency ωα mode
from the analysis at the outset by setting up the problem such that the center of mass of the
system is fixed. The center of mass (along the x axis) of the system is
1 X 1
xcm = P mi xi = [m(x1 + x3 ) + Mx2 ] . (4.235)
mi 2m + M
By choosing our coordinate system such that xcm = 0, we can eliminate one of the coordi-
nates, say x2 , by writing
m
x2 = − (x1 + x3 ) . (4.236)
M
Substituting this back into Eqs. (4.215)-(4.217) yields two independent equations of mo-
tion, which can be solved in the usual fashion with the result that only the normal frequen-
cies corresponding to ωβ and ωγ appear as solutions. We leave this as an exercise at the
end of the chapter.
t
Fig. 4.12 Displacements for normal modes of CO2 : (a) longitudinal modes α, β, and γ, along the
x direction, (b) transverse mode δ along the y direction. The α mode is not a
vibrational mode but corresponds to pure translation. Other pure translational and
rotational modes are not shown.
t
4.3.2 Transverse modes
Thus far, we have only allowed for motion along a line colinear with the equilibrium po-
sitions of the atoms, i.e. along the x axis. The modes thus obtained are called longitudinal
modes. There also exist transverse modes, that is, modes perpendicular to the x axis, which
we can study by allowing for motion in the y and z directions. However, as we have seen
from our analysis of the longitudinal modes, the number of modes we obtain can depend
on how we set up the problem. Before diving into another calculation, perhaps it is best
to think a bit. In particular, it would be useful to know how many transverse modes we
are looking for. On quite general grounds, we know that the number of modes is equal to
the number of degrees of freedom. Since each particle has three degrees of freedom and
there are three particles, there are a total of 3 × 3 = 9 degrees of freedom. Let’s count the
different types. There are 3 translational degrees of freedom in which the whole molecule
moves in a straight line along the x, y, or z direction. There are 2 rotational degrees of free-
dom: the molecule can rotate about the y or z axis (but not about the x axis). In the previous
section, we found 2 longitudinal degrees of freedom. That leaves us with 9 − 3 − 2 − 2 = 2
transverse degrees of freedom. Based on the linear symmetry of the CO2 molecule, we
know that for any mode we find that involves motion only in the y direction, there must be
a corresponding completely independent mode oscillating at the same frequency in the z
direction. When two modes such as these have the same frequency, we say they are degen-
erate.4 These 2 degenerate transverse modes account for the last 2 modes of the system.
Since the two modes are degenerate, we have only one more mode to find.
Transverse y modes
We begin by determining the center of mass in the y direction and setting it equal to zero
in order to preclude the sort of trivial translational mode we encountered in the previous
section:
1 X 1
ycm = P mi yi = [m(y1 + y3 ) + My2 ] . (4.237)
mi 2m + M
Setting ycm = 0, we obtain

m
y2 = − (y1 + y3 ) . (4.238)
M
The symmetry of having two identical masses on either side of the central mass suggests
normal modes where either: (1) the two outer masses move in opposite directions along
the y axis and the central mass remains still so that ycm = 0, or (2) the two outer masses
move in the same direction along the y axis while the central mass moves in the opposite
direction so that ycm = 0. Indeed this is what we found for the longitudinal modes. Mode
(1) corresponds to the beginning of a rotational mode, which we have already counted and
excluded from this analysis. That leaves us with mode (2).
4 Degenerate modes are, in spite of their name, are no more reprobate than any other mode.
114 Normal Modes
t
The equation of motion for the three masses are
F1 = −κlθ = mÿ1 (4.239)

F2 = 2κlθ = M ÿ2 (4.240)
F3 = −κlθ = mÿ3 , (4.241)
where κ is the torsional spring constant for bending. Note that there is no obvious rela-
tionship between the torsional bending constant κ and the longitudinal spring constant k
introduced previously. The angle θ is related to the transverse displacements by
(y1 − y2 ) + (y3 − y2 )
θ= , (4.242)
l
in the limit of small displacements where yi l. For the mode we seek, y3 = y1 , yielding
2(y1 − y2 )
θ' . (4.243)
l
Using Eq. (4.238) and noting once again that y3 = y1 , we obtain y2 = −(2m/M)y1 , which
we can use to eliminate y1 from our expression for θ. This gives
2y2 M
θ'− 1+ . (4.244)
l 2m
Substituting this expression for θ into Eq. (4.240) gives the equation of motion subject to
the constraints that ycm = 0 and y1 = y3 :
M
M ÿ2 + 4κ 1 + y2 = 0 . (4.245)
2m
This is just the usual equation for a simple harmonic oscillator with an oscillation frequency
r
4κ M
ωδ = 1+ (4.246)
M 2m
The normalized eigenvector for the δ mode is
 1 
 
1
aδ = p (4.247)
 
−2m/M 
2 + 4(m/M) )
2 
1

Exercise 4.3.1 Starting from Eqs. (4.239)-(4.241) and Eq. (4.242):

(a) Write down the mass and stiffness matrices m and k for the transverse modes
in terms of the coordinates (y1 , y2 , y3 ). You should find that both m and k are
symmetric.
(b) Calculate the matrix D = m −1k (which in this case is not symmetric) and verify
by direct substitution that (1, −2m/M, 1) is an (unnormalized) eigenvector with
a normal frequency given by Eq. (4.246).
(c) Verify that (1, 1, 1) and (1, 0, −1) are unnormalized eigenvectors with zero fre-
quency. Explain what each of these modes means physically. Does either mode
pose any problem for the yi 1 approximation we made?
t
(d) Extra: Show that (2, 1, 0) is also an eigenvector with an eigenvalue of zero.
This mode cannot correspond to another independent normal mode because we
have already found the three independent transverse y modes that are allowed
for a system with three degrees of freedom. Show that the (2, 1, 0) is not an
independent mode.
Degeneracy: transverse y and z modes

In the previous section we found a single transverse normal mode in the y direction with
a nonzero frequency. Of course, that mode has an identical counterpart oscillating in the z
direction. The frequency of both modes is given by Eq. (4.246). The transverse oscillations
in the y direction are given by
y(t) = êy aδ (aδy cos ωδ t + Bδy sin ωδ t) , (4.248)
where êy is a unit vector in the y direction, ωδ is given by Eq. (4.246), aδ is given by Eq.
(4.247), and the constants aδy and Bδy are determined by the initial conditions. Similarly,
transverse oscillations in the z direction are given by
z(t) = êz aδ (aδz cos ωδ t + Bδz sin ωδ t) . (4.249)
where the constants aδz and Bδz are determined by the initial conditions.
Figure 4.13(a) shows the geometry for the two degenerate transverse modes. The direc-
tions y and z are arbitrary, of course, as long as we choose them to be perpendicular to the
molecular axis of the CO2 molecule. The oscillation frequency of these modes corresponds
to the oscillation frequency of infrared light with a wavelength of 15 µm. If we illuminated
a container of CO2 with 15 µm infrared radiation that was linearly polarized along a certain
direction, then it would excite the CO2 molecules to oscillate in that direction in a trans-
verse mode. In that case, it would be convenient to choose the y direction to correspond to
only one of the two degenerate modes. Otherwise, we would have to describe it as a linear
combination, the vector sum, of the two degenerate modes. There is nothing wrong with
doing this, other than making the problem a bit more complicated than necessary.
Suppose we were exciting the vibrational modes of CO2 with circularly polarized light.
Then we would expect that the atoms in the CO2 molecules would respond by vibrating in
one of the circularly polarized vibrational modes depicted in Fig. 4.13(b)-(c). Such modes
t
Fig. 4.13 Transverse modes of a carbon dioxide molecule: (a) y and z linear polarized modes,
(b) right circularly polarized, (c) left circularly polarized.
116 Normal Modes
t
can be readily created from linear combinations of Eqs. (4.248) and (4.249) by making the
two modes oscillate 90◦ out of phase with each other. For example, the left and right circu-
larly polarized modes shown in Fig. 4.13(b)-(c) can be realized by setting aδy = ±Bδz = 1
and Bδy = aδz = 0. In fact, any linear combination of the transverse modes is valid because
of the underlying linearity of the equations of motion.
If we are working with the circularly polarized modes, then it is a bit cumbersome to
describe them as phased superpositions of linearly polarized modes. It would be more
convenient if we could express them mathematically in a direct way. To facilitate this, we
first rewrite the linearly polarized modes using complex notation. Recalling that in so doing
we ultimately keep only the real part of our complex expressions, Eqs. (4.248) and (4.249)
become
y(t) = êy aδ aδy eiωδ t (4.250)

iωδ t
z(t) = êz aδ aδz e , (4.251)
where aδy and aδy are in general complex, and thus are able to carry the phase information
that is contained in the sine and cosine functions in Eqs. (4.248) and (4.249).
It is a simple matter to create circularly polarized modes from Eqs. (4.250) and (4.251)
above. Recalling that we keep only the real part, if we choose aδz = e±iπ/2 aδy = ±iaδy , the
we advance or delay the transverse z mode relative to the y mode, producing right and left
circularly polarized states, which we can write as
1
L(t) = √ aδ aδy eiωδ t êy + iêz

(4.252)
2
1
R(t) = √ aδ aδy eiωδ t êy − iêz ,

(4.253)
2
where the normalization factor of 2−1/2 assures that the column matrix has a length of unity.
4.4 Damping and normal modes
The systems of coupled oscillators that we have examined thus far have been undamped.
Real systems exhibit damping, of course, so we proceed here to see what happens when
damping in included. The matrix formulation really comes in handy in this situation, as it
can teach us some general features about damped coupled oscillators that otherwise might
be difficult to figure out. To this end, we write down the equations of motion of a system
of damped coupled oscillators in matrix form:
m ẍ = −k x − Γ ẋ . (4.254)
We proceed as before (see Eq. (4.80)) and substitute a trial solution of the form
x = a e−iωt , (4.255)
117 Damping and normal modes
t
which yields
−ω2m a = −k a + iωΓ a . (4.256)
Left multiplying both sides by m −1 and rearranging terms so that the frequency ears only
on the right hand side gives
m −1k a = ω2 a + iω m −1Γ a . (4.257)
Equation (4.257) is no longer an eigenvalue equation, which means that the normal mode
analysis we have developed cannot be used to find solutions. Nevertheless, Eq. (4.257) still
represents a system of N equations with N unknowns where N is the number of degrees
of freedom of the system of coupled oscillators. While it is possible to find the solutions
to these equations, the analysis can be tedious. So we will not attempt to find general
solutions to this equation. Instead we consider a special case for which Eq. (4.257) reduces
to an eigenvalue problem. Solutions for the more general case qualitatively resemble those
to the special case we are considering, and thus enable us to explore the basic physics of
the problem.
The special case we consider is the case when the damping matrix is proportional to the
mass matrix, that is, when
Γ ∝m . (4.258)
In this case, m −1Γ ∝ I so that Eq. (4.257) can be rewritten as an eigenvalue equation,
D a = λa , (4.259)
where D = m −1k and λ is a scalar, perhaps complex, function of ω. More generally, in

transforming Eq. (4.256), we would want to obtain a symmetric dynamical matrix, meaning
we would want to follow the procedure we developed in §4.2.5. Let’s work through an
example.
Example 4.3 Consider the coupled pendulum we introduced at the beginning of this
chapter, but for the case where each pendulum experiences a damping force proportional
to its velocity. Assume the damping constants are the same for the two pendulums. Find
the normal modes, their frequencies, and the damping for each mode.
Solution
We start by writing down the equations of motion for the masses of the two pendulums,
including a damping term:
x1
F1 = −mg + k(x2 − x1 ) − Γ ẋ1 = m ẍ1 (4.260)
l
x2
F2 = −mg − k(x2 − x1 ) − Γ ẋ2 = m ẍ2 , (4.261)
l
where we have included identical damping terms for each pendulum, and have used the
same notation and small angle approximations used previously to for the undamped case
118 Normal Modes
t
treated in §4.1.1. Collecting terms together and rearranging gives
mg
m ẍ1 = − + k x1 + kx2 − Γ ẋ1 (4.262)
l
mg
m ẍ2 = kx1 − + k x2 − Γ ẋ2 , (4.263)
l
These two equations can be rewritten in matrix form following Eq. (4.254), with the mass
and stiffness matrices given by Eq. (4.183) with M = m and the damping matrix given by
Γ 0 m 0
! !
Γ= =γ = γm , (4.264)
0 Γ 0 m
where γ ≡ Γ/m. Because Γ = γm , the matrix m −1Γ appearing in Eq. (4.257) is given by
m −1Γ = m −1 γ m = γI . In this case, Eq. (4.257) reduces to
D a = ω2 a + iγω a (4.265)
= λa . (4.266)
where D = m −1k and λ = ω2 + iγω. In this case D is symmetric because the two masses are
equal, so we can proceed using D = m −1k . Because Γ ∝ m the equations of motion can be
written in the standard form of an eigenvalue equation. Moreover, because m and k are the
same as for the undamped case, D is unchanged and the eigenvalue equation is formally the
same as it was for the undamped case. Thus, the eigenvalues are the same, namely λ = ω2p
and λ = ω2p + 2ω2b . The only difference is that for the damped case λ = ω2 + iγω whereas
for the undamped case λ = ω2 . Therefore the frequencies are given by solutions to
ω2 + iγω − ω2p = 0 (4.267)
ω2 + iγω − ω2p + 2ω2s = 0 .

(4.268)
From the quadratic formula we obtain
r γ 2
iγ
ωα = − ± ω2p − (4.269)
2 2
r γ 2
iγ
ωβ = − ± ω2p + 2ω2s − . (4.270)
2 2
These solutions are the same as those obtained in §2.1 for a single damped oscillator where
the normal mode frequency for each mode replaces the natural frequency ω0 of the pendu-
lum. Each mode of the coupled pendulums, therefore, exhibits the same kind of behavior
as the simple pendulum, with the same underdamped, overdamped, and critically damped
cases. The only difference is that there are two “natural frequencies” for the coupled pen-
dulum, namely ωα and ωβ , whereas there is only one for the simple pendulum, namely
ω0 . This leads to somewhat more complex damping for the coupled pendulums. Although
there are two normal frequencies, Eqs. (4.269) and (4.270) have only a single damping rate
γ/2, which is common to both modes. This is a consequence of our assumption the Γ ∝ m
and the fact that we took both pendulums to have the same mass.
Finally we note that the normal modes for the damped pendulum are exactly the same as
those for the undamped pendulum: there is a low-frequency α mode where the pendulums
119 Forced oscillations
t
oscillate in phase and a high-frequency β mode where the pendulums oscillate out of phase.
The difference in this case is simply that the oscillations die out with time.
4.5 Forced oscillations
Coupled oscillators are often driven by some external force, just as simple oscillators are,
and not surprisingly, this leads to the phenomena of resonance, with some interesting new
twists. Intuitively, we might expect that there will be a resonance at each normal frequency,
and indeed, this is what happens, albeit with a few caveats. Here we embark on a brief
treatment of forced oscillations and resonance in coupled oscillators.
To examine forced oscillations in coupled oscillators, we return once again to the prob-
lem of two coupled pendulums. We also include damping, but in such a way that the damp-
ing matrix is proportional to the mass matrix, just as we did in the previous section. This
keeps the analysis fairly simple and illustrates the basic phenomena. We will also examine
the limit of negligible damping.
4.5.1 Steady-state response
To start, we include sinusoidal forcing only for pendulum 1, the left pendulum, in Fig. 4.1.
For this case, we assume no forcing of the second pendulum. Starting from Eqs. (4.262)
and (4.263), the equations of motion become
mg
m ẍ1 = − + k x1 + kx2 − Γ ẋ1 + F0 e−iωt (4.271)
l
mg
m ẍ2 = kx1 − + k x2 − Γ ẋ2 , (4.272)
l
where F0 e−iωt is a sinusoidal forcing term.
Before solving the problem, it’s useful to think about the underlying physics. We know
that there are two normal modes from our analysis of the undamped and damped problem:
the slow α mode with the normal coordinate qα = √12 (x1 + x2 ), and the fast β mode with
the normal coordinate qβ = √12 (x2 − x1 ), which are just the normalized versions of Eqs.
(4.22) and (4.23), respectively. This means that the position x1 of the driven pendulum is a
linear combination of qα and qβ : that is, x1 = √12 (qα + qβ ). Thus, when we drive pendulum
1 with an external force, we are exciting both qα and qβ : that is, both the slow α mode
and the fast β mode. For a given drive frequency, we would not expect the two modes to
respond equally, however, as their respective responses should depend on how close the
driving frequency ω is to the normal frequencies ωα and ωβ of the two modes. The closer
the drive frequency is to a given normal frequency, the more that mode would be expected
to respond. With these ideas in mind, let’s turn to the solution of the problem.
The general approach to finding the mathematical solution to such a problem is to first
120 Normal Modes
t
find the normal modes for the problem without external driving using the matrix formula-
tion. Thus, we start by restating the problem in matrix form. Following our development
in §4.2.1 and §4.4, the equation of motion is
m ẍ = −k x − Γ ẋ + F0 e−iωt . (4.273)
where m and k are given by Eq. (4.183), and Γ is given by Eq. (4.264). The column matrix
(or vector) F0 is given by
F0
!
F0 = , (4.274)
0
where the top element is finite and the bottom element is zero because only pendulum 1 is
driven.
To solve this problem, we return to the formalism developed in §4.2.2. We will also de-
velop the formalism allowing for the possibility of unequal masses as developed in §4.2.5.
The first step is to transform to a basis in which the dynamical matrix is diagonal. Following
the procedure of §4.2.5 of left multiplying both sides by m −1/2 and inserting m −1/2 m 1/2 = I
in strategic places, Eq. (4.273) becomes
m −1/2m ẍ = −m −1/2k m −1/2 m 1/2 x − m −1/2Γ m −1/2 m 1/2 ẋ + m −1/2 F0 e−iωt . (4.275)

which reduces to
ÿ = −D y − m −1/2 γ m m −1/2 ẏ + m −1/2 F0 e−iωt (4.276)

= −D y − γ ẏ + m −1/2
F0 e−iωt
, (4.277)
where y ≡ m 1/2 x. We now transform to the ỹ basis using Eq. (4.107)
ỹ = U y , (4.278)
where U T is the matrix whose columns are the normalized eigenvectors of D . Multiplying
Eq. (4.277) by U
U ÿ = −UDU T U y − γ U ẏ + Um −1/2 F0 e−iωt (4.279)

ỹ¨ = −D̃
Dỹ − γ ỹ˙ + Um −1/2 F0 e−iωt , (4.280)
D ≡ UDU T and we have inserted U T U = I , noting that U −1 = U T . If there is no

where D̃
damping or forcing, Eq. (4.280) becomes
ỹ¨ = −D̃
Dỹ . (4.281)
We saw already in Eq. (??) that in the tilde basis, the equations of motion are decoupled.
This means that the matrix D̃D is diagonal with diagonal entries that are the eigenvalues
of the system, in this case ω21 and ω22 . As the equations are decoupled, we write out Eq.
(4.280) in component form,
ỹ¨ 1 + γỹ˙ 1 + ω21 ỹ1 = F̃01 e−iωt (4.282)

ỹ¨ 2 + γỹ˙ 2 + ω22 ỹ2 = F̃02 e−iωt , (4.283)
121 Forced oscillations
t
where
F̃0 = Um −1/2 F0 . (4.284)
These are the same as the equation of motion for a single forced oscillator that we encoun-
tered in §3.1.1 (see Eq. (3.2)).
The eigenvectors for equal-mass coupled pendulums are
sin θ 1 1 − sin θ 1 −1
! ! ! !
ẽ1 = = √ , ẽ2 = = √ , (4.285)
cos θ 2 1 cos θ 2 1
where θ = 45◦ . Since the columns of U T are the eigenvectors of the system
√
1 1 1 1/ m 0 F0
! ! !
F̃0 = Um −1/2 F0 = √ √ (4.286)
2 −1 1 0 1/ m 0
1 F0 1
!
= √ √ , (4.287)
2 m −1
√
which means that F̃01 = F̃02 = F0 / 2m.
For the example here where the masses are equal, we can trivially transform between
the
Because the form of the differential equations describing the normal modes, Eqs. (4.282)
and (4.283), is the same as the differential equation describing a single forced oscillator,
Eq. (3.4), we can adapt the solution to Eq. (3.4) directly to Eqs. (4.282) and (4.283). Thus,
using the solution expressed in Eq. (3.15) as a guide together with Eqs. (3.12) and (3.13),
the solutions for the normal modes qα and qβ are given by
qα (t) = aα (ω) cos(ωt − φα ) (4.288)

qβ (t) = −aβ (ω) cos(ωt − φβ ) , (4.289)
where
√1 F 0 /m
γω
!
2
aα (ω) = p , φα = tan−1 , (4.290)
(ω2α − ω2 )2 + γ2 ω2 ω2α − ω2
√1 F 0 /m
 
2
 γω 
aβ (ω) = q , φβ = tan  2
−1   . (4.291)
(ω2β − ω2 )2 + γ2 ω2 ωβ − ω2 
The responses of the normal modes given by Eqs. (4.288)-(4.291) are simple enough as
they mimic exactly the results we obtained for a single resonator. It is interesting to look at
the response of the individual masses, however, as they are a little less evident. For the case
of the coupled pendulums considered here, Eqs. (??) and (??) can be inverted to express
the coordinates of the individual masses in terms of the normal coordinates
1
x1 = √ (qα − qβ ) (4.292)
2
1
x2 = √ (qα + qβ ) . (4.293)
2
122 Normal Modes
t
A1
4
3
2
1
0 t
2 3 4
−1
−2
−3
A2
0 t
2 3 4
−1
−2
−3
t
−4
Fig. 4.14 Resonance curves for two weakly coupled pendulums with forcing: Amplitudes of
masses 1 and 2 for the slow and fast modes as a function of forcing frequency.
These two equations together with Eqs. (4.288)-(4.291) give the steady state response of
the two masses, which are plotted in Fig. 4.14.
4.6 Summary of important points of Chapter 4
• The motion of a system of coupled linear oscillators can be described in terms of its
normal modes.
• The number of normal modes N of a system is equal to the number of degrees of free-
dom.
• For each normal mode, there is a normal coordinate qα that is a linear combination of
the original coordinates used to describe the system.
• When the coordinates of all of the particles in the system are expressed in terms of the
normal coordinates, the system of N second order differential equations decouples into
N uncoupled differential equations of the form q̈α + ω2α q = 0, where ωα is the (normal)
frequency of that mode.
• For each normal mode there is one normal frequency.
• Two or more normal modes are said to be degenerate if they have the same normal
frequency. Degenerate normal modes generally are the result of some kind of symmetry
of the system allowing the system to oscillate in essentially the same way along two or
more equivalent directions. For example, the transverse modes of the CO2 molecule are
degenerate because oscillations along the transverse directions (the y and z directions in
Fig. 4.13)) are identical.
123 Problems
t
Problems
4.1 Two simple pendulums of length 0.300 m and and mass 0.950 kg are coupled by
attaching a light horizontal spring of spring constant k = 1.50 N/m to the masses.
(a) Determine the frequencies of the two normal modes.
(b) One of the pendulums is held at a small distance away from its equilibrium posi-
tion while the other pendulum is held at its equilibrium position. The two pendu-
lums are then released simultaneously. Show that after a time of approximately
12 s the amplitude of the first pendulum will become equal to zero momentarily.
(assume g = 9.81 m/s2 ).
4.2 (French 5.5) Two identical undamped oscillators, A and B, each of mass m and
natural (angular) frequency ω0 , are coupled in such a way that the coupling force
exerted on A by B is αm(d2 xB /dt2 ), and the coupling force exerted on B by A is
αm(d2 xA /dt2 ), where α is the coupling constant of magnitude less that 1. Describe
the normal modes (i.e., specify the relative amplitudes of the displacements of mass
ω20
A and mass B) of the coupled system and find their frequencies. [Ans: ω2± = (1±α)
with respective amplitude rations AB /AA = ∓1.]
4.3 (French 5-6) Two equal masses on an effectively frictionless horizontal air track are
held between rigid supports by three identical springs, as shown in the figure below.
The displacements from equilibrium along the line of the springs are described by the
coordinates xA and xB , as shown. If either of the masses is clamped in its equilibrium
position, the period T = 2π/ω for one complete vibration of the other mass (attached
to two springs) is 3 sec.
(a) If both masses are free (but still attached to the springs), what are the periods of
the two normal√ modes of√the system? Sketch a graphs of xA (t) and xB (t) in each
mode. [Ans: 6 s and 3 2 s.]
(b) At t = 0, mass A is at its normal resting position and mass B is pulled aside a
distance of 5 cm. The masses are released from rest at this instant.
1. Write an equation for the subsequent displacement for each mass as a function
of time.
2. What is the period (in seconds) for the transfer of from B to A and back.
cycle is the situation at time t = 0 exactly reproduced? Explain.
After one √
3
[Ans: 2 ( 3 + 1)]
√
4.4 Two masses m are suspended in series from springs with the same force constant k
as shown in the figure below:
124 Normal Modes
t
(a) Find the normal frequencies and normal coordinates of this system. Hint: Define
your coordinates for the two masses with respect to the static equilibrium posi-
tions of each mass. In this way, you should find that constant terms proportional
to mg disappear from the equations of motion. [Ans:
r
1 √ 1 √
ω1 = (3 − 5) ωa ẽ1 = q (2, 1 + 5)
2 √
10 + 2 5
r
1 √ 1 √
ω2 = (3 + 5) ω1 ẽ2 = q (−(1 + 5), 2),
2 √
10 + 2 5
√
where ωa = k/m.]
(b) At t = 0, both masses are at their respective equilibrium positions each mov-
ing with the same velocity 30 . Find expressions for the motion of each mass as
function of time given these initial conditions. Ans:
 
30 /ωa 2 4
 
x1 (t) = √ sin ω t + sin ω t
 
√  q √ 1
√ 2 
2(5 + 5)  3 − 5
q 
3+ 5
 
 
√ √
30 /ωa  4(2 + 5) 2( 5 − 1)
 
x2 (t) = √ sin ω t sin ω t

√  q 1 − 2 
2(5 + 5)  3 − 5 √ q √ 
3+ 5
 
4.5 Consider a pair of pendulums, one twice as long as the other coupled to each by a
spring. Each pendulum consists of a massless stiff rod with a mass m at its end. The
spring is attached halfway along the length of the longer pendulum and to the end of
the other pendulum as shown in the figure below. Assume that θ1 = θ2 = 0 and that
the spring is unstretched at equilibrium.
(a) Write down the coupled equations of motion for the two pendulums, using the
angles θ1 and θ2 as the dynamical variables and assuming that θ1 1 and θ2
1. Write the coupled equations of motion in angular form using torque τ ≡ rF,
that is, in the form τ = I θ̈, where I is the moment of inertia. Show that the
equations of motion can be written in matrix form
m θ̈ = −k θ , (4.294)
125 Problems
t
where
ml2 0 mgl + kl2 −kl2 θ1
! ! !
m= , k= , θ= . (4.295)
0 4ml2 −kl2 kl2 + 2mgl θ2
(b) Show that the equation of motion can be rewritten in the standard form θ̈ = −D 0 θ
with
ω2 + ω2 −ω2s
!
D 0 = m −1k = p 2 2s , (4.296)
−ρ ω s ρω2p + ρ2 ω2s
where ω2p = g/l and ω2s = k/m. What is the numerical value of ρ?
(c) Find the symmetric form of the dynamical matrix given by D = m −1/2k m −1/2 .
√
(d) For the remaining parts of this problem set ω p = 1 and ω s = 2. Find the
normal frequencies of the system.
(e) Find the normalized eigenvectors of D . Make a plot of the eigenvectors in the
y basis defined by y = m 1/2 θ. Writing the normalized eigenvectors as in Eq.
(4.103), show that the eigenvectors are rotated by an angle of 67.5◦ in the y
basis.
(f) Find equations for the trajectories of each pendulum as a function of time for the
initial conditions where θ1 = θ2 = 0, θ̇1 = Ω > 0, and θ̇2 = 0. Make plots of θ1 (t)
and θ2 (t).
4.6 In this problem you analyze the normal modes of the Wilberforce pendulum, which
consists of a body of mass m and moment of inertia I hanging from a spring, as shown
in Fig. 4.15. When a Wilberforce pendulum is set into motion, say by stretching
the spring a small distance from its equilibrium position, the pendulum begins to
oscillate up and down. As time passes, the pendant mass slowly begins to twist about
its vertical axis while the amplitude of the vertical oscillations becomes smaller due
to a weak coupling between extension and torsion of the spring. As time continues to
pass the system periodically exchanges energy between the extensional and torsional
motion.
The coupling between extension and torsion can be understood as follows. When
the mass on the spring is displaced vertically from its equilibrium position, the radius
of the coil of the spring changes slightly, causing a small bit of torsion in the spring
and changing the equilibrium twist of the spring. Thus, the potential energy of the
126 Normal Modes
t
spring is given by
1 2 1
U(z, ϕ) = kz + η[ϕ − ϕ0 (z)]2 , (4.297)
2 2
where ϕ0 (z) gives the angular shift in the equilibrium twist when the spring is dis-
placed a vertical distance z from its equilibrium position. Because the displacements
are small, we can take the shift to be linear in z
ϕ0 (z) = bz . (4.298)
Substituting this into Eq. (4.297) gives

1 1
U(z, ϕ) = (k + ηb2 )z2 + ηϕ2 − ηbzϕ . (4.299)
2 2
Typically, the coupling between extensional and torsional motion is weak, meaning
that ηb2 k.
(a) The vertical force exerted by the pendulum on the mass is given by Fz = −∂U/∂z;
the torque (torsional force) is given by τϕ = −∂U/∂ϕ. Use these equations and
the expression for the potential energy above to obain the following coupled set
of equations of motion for the Wilberforce pendulum:
d2 z
m = −Kz + ϕ (4.300)
dt2
d2 ϕ
I 2 = z − cϕ , (4.301)
dt
where the constants K, , and c are various combinations of k, η, and b that you
should determine.
(b) What are the natural frequencies ωz and ωϕ of the translational and rotational
oscillations, respectively, in the absence of the coupling terms in Eqs. (4.300)
and (4.301), that is, when b = 0.
t
Fig. 4.15 Wilberforce pendulum twisted an angle ϕ − ϕ0 (z) past its equilibrium position ϕ0 (z) for
a vertical displacement of z.
127 Problems
t
(c) Show that the two normal frequencies are given by
s
2 K η K η 2 2
ωα = + − − + (4.302)
2m 2I 2m 2I mI
s
K η K η 2 2
ω2β = + + − + (4.303)
2m 2I 2m 2I mI
(d) In our class demonstration, we observed the periodic exchange of energy be-
tween vertical and torsional motion in the Wilberforce pendulum, similar to that
plotted in Fig. 4.4. One can show that beats occur like those shown in Fig. 4.4
when ωz = ωϕ . The vertical oscillation period was about 0.8 s and the beat
period—that is the time it took the pendulum to go through one full cycle of no
torsional oscillations through a maximum in torsional oscillations and back to no
torsional oscillations—was about 20 s. Given that the mass was a steel cylinder
about 2 cm high and about 4 cm in diameter, find the approximate values of the
mass m, the moment of inertia I, the spring constant k, and the torsional spring
constant η. The value of b is measured to be about 2.5 rad/m.
4.7 Consider a system exactly like the the pendulum swinging from a sliding mass pre-
sented in §?? but without a spring attached to the mass M. The mass M is still free
to slide without friction on the surface supporting it. Find the normal modes and
frequencies of the system and describe the motion of each mode physically.
4.8 Consider three equally spaced masses, M, m, and m, connected by springs and con-
strained to move on a circle of radius R that is fixed in space, as shown in the figure
below. All three springs have the same spring constant k. The dashed spokes show
particle positions with θ1 = θ2 = θ3 = 0 corresponding to unstretched springs; dotted
lines show positive (clockwise) displacements.
(a) Write down the equations of motion and cast the problem in terms of matrices
following the formalism introduced in §4.2.1. Using the matrix formalism, find
all the normal modes of the system including the normal frequencies and the nor-
malized eigenvectors (normal coordinates). Make a sketch of the normal modes
and briefly discuss the motion associated with each one.
128 Normal Modes
t
(b) Starting from your solution to part (a), find the equations for the motion of the
three masses for the following set of initial conditions: θ1 = 0, θ̇1 = Ω, θ2 = 0,
θ̇2 = −Ω, θ3 = 0, θ̇3 = 0,
4.9 In §4.3.1 we found three longitudinal modes of the carbon dioxide molecule, one of
which was the zero-frequency mode corresponding to simple translation of the entire
molecule at constant velocity. Reformulate the problem as one with only two degrees
of freedom by fixing the center of mass to be zero. That is, find the center of mass in
terms of x1 , x2 , x3 , m, and M. By insisting that the center of mass remain fixed, say
at a value of zero, eliminate one of the three variables x1 , x2 , or x3 , and reformulate
the problem with only two degrees of freedom. Then find the normal frequencies and
eigenvectors using the matrix formulation.
4.10 Consider the two masses hanging from springs that you considered in Problem 4.4.
Suppose the bottom mass if driven with the sinusoidal force F(t) = F0 cos ωt and
that the motion of each mass is damped with a damping force Fd = Γ3, where 3 is
the velocity of the particular mass.
(a) Using the results from your solution to Problem 4.4 (or the one posted by your
TA), find the steady state response of the system to the external drive. In particu-
lar, find equations for the motion of the system in terms of the normal coordinates
and for the motion of the individual masses.
(b) What does the motion of the two masses look like when the drive frequency is
equal to that of the slow normal mode? the fast normal mode?
4.11 (a) In §4.1.1, we found the normal coordinates for the two pendulums coupled by
a spring. Using the matrix formalism developed in §4.2.2 show that the ma-
trix U −1 that transforms from the (x1 , x2 ) coordinates to the ( x̃1 , x̃2 ) normal co-
ordinates is a rotation through 45◦ . That is, show that U −1 is the conventional
two-dimensional rotation matrix (for a specific angle). For information about the
rotation matrix, see http://en.wikipedia.org/wiki/Rotation_matrix.
(b) In Appendix B, we solve for the normal modes of a double pendulum using the
Lagrangian formalism to obtain the equation of motion. Using the results of that
analysis, find the transformation matrix U −1 that transforms from the (x1 , x2 ) co-
ordinates to the (qα , qβ ) for the double pendulum. Can that transformation be
described as a rotation between the the (θ, φ) coordinates and the ( x̃1 , x̃2 ) normal
coordinates? That is, is the transformation matrix U −1 a rotation matrix for some
angle? Using the dot product, find the angle between ( x̃1 , x̃2 ) for the double pen-
dulum. Finally, find the normal coordinates of the double pendulum and show
by substituting them into Eqs. (B.35) and (B.36), the small-angle equation of
motion, that the equations decouple.
5 Waves of oscillating particles
Most of the waves we encounter in everyday life involve the excitation of some medium:
waves on the strings of a guitar or violin, the undulations of the surface of a pool of water
or of a drum, or sound waves in air, water, or through a solid wall. Such waves involve
particles moving up and down or back and forth, and are properly viewed as normal modes
such as those introduced in Chapter 4. There, we studied the normal modes of systems
consisting of a few particles, two, or at most three. Here, we consider waves involving a
much larger number of particles. As we shall see, the fact that particles are separated a
finite distance from each other imposes a lower cutoff on the length scale on which waves
can exist, which has important physical consequences. When particles are neatly arranged
on a lattice, this leads to phenomena such as band gaps, that is, the absence of propagating
waves over certain frequency ranges. When the arrangement of particles is disordered, or
the interactions between particles vary randomly in strength, the normal modes can become
highly localized in space, involving only a subset of nearby particles within the system. In
fact, the phenomena involving waves on arrays of interacting particles is almost limitless,
and are sufficiently subtle and important to remain the focus of a great deal of current
research.
In this chapter, we begin our inquiry into the phenomena of waves of oscillating particles
by considering a deceptively simple but revealing example: point masses on a massless
string. We begin with N identical masses evenly spaced along a string, and then explore
the consequences of having N pairs of two different masses evenly spaced along a string,
which leads to the appearance of a band gap. We then explore the consequences of other
arrangements of particles of different masses along a string, in particular, quasicrystalline
and disordered arrangements, both topics of current research. All of these examples are
drawn from systems that are essentially one dimensional, meaning that the particles are
always arranged along a line. This assumption simplifies the math and illustrates much of
the essential physics.
5.1 Identical point masses on a massless string
We start by considering the problem of N masses tethered to each other by a string stretched
between two fixed walls, as shown in Fig. 5.1. The masses are evenly distributed along the
string, with each mass a distance l from its nearest neighbor. At equilibrium, the masses and
string form a straight line. We shall consider the oscillations that occur when the masses
129
t
are displaced perpendicular (transverse) to the horizontal line connecting the masses at
equilibrium, as shown in Fig. 5.1. We index the masses by the letter s. The transverse
displacement of the sth mass is given by y s . The ends of the string are fixed to walls and do
not move so that y0 = yN+1 = 0. We consider the case where the transverse displacements
are small meaning that the angle made with he horizontal by the sth string segment θ s 1
(see Fig. 5.1(b)). Taking the tension T along the string to be constant, which is valid for
small transverse displacements, the equation of motion the sth mass is
Fys = T sin θ s − T sin θ s−1 = mÿ s . (5.1)
For small transverse displacements, we can make the small angle approximation that sin θ s '
(y s+1 − y s )/l and sin θ s−1 ' (y s − y s−1 )/l, yielding the following equation of motion
T
mÿ s = (y s−1 − 2y s + y s+1 ) . (5.2)
l
Equation (5.2) actually represents N coupled equations of motion, one for each particle on
the string. Thus, as the index s runs from 1 to N, we obtain:
T
mÿ1 = (−2y1 + y2 ) (5.3a)
l
T
mÿ2 = (y1 − 2y2 + y3 ) (5.3b)
l
..
.
T
mÿN−1 = (yN−2 − 2yN−1 + yN ) (5.3c)
l
T
mÿN = (yN−1 − 2yN ) . (5.3d)
l
Each equation has three dynamical variables, those involving a particle and its two nearest
neighbors, which couples each equation to the equation before and after it. The only excep-
tions are the first and last equations, which have only two terms because of the boundary
conditions y0 = yN+1 = 0 discussed above.
We can recast the set of equations in Eq. (5.3) as a matrix equation with the same form
that we encountered in Chapter 4:
m ÿ = −k y , (5.4)
t
Fig. 5.1 N masses on a string under tension T . At equilibrium, the masses form a straight
horizontal line. (a) Transverse displacements {y s } of N particles of mass m separated
by a horizontal distance l on a string. (b) Close-up showing forces on the sth mass m
tethered to a string under tension T .
131 Identical point masses on a massless string
t
where
m 0  1 0  y1 

     
m 1  y2 
 
  
  
   
m = 
 m  = m 


 

 1  , y =  y3 


  
.. ..  . 

.  .. 
 
.

   
    
0 m 0 1 yN
   
 2T /l −T /l 0   2 −1 0 

   
−T /l 2T /l −T /l −1 2 −1
   
 
 ..  T  .. 
k =  −T /l 2T /l .  =  −1 2 .  .
 l  
.. .. .. ..

. . . .
  
−T /l
  
  −1
0 −T /l 2T /l 0 2
   
−1
With these definitions, performing the matrix multiplications implied by Eq. (5.4) yields
Eqs. (5.3a)-(5.3d). Note that the stiffness matrix k is tridiagonal, that is, it has non-zero
terms only along the diagonal and the two adjacent entries in the same row. That these are
the only non-zero terms reflects the physical fact that each mass is coupled through the
string only to its nearest neighbor. By contrast, if the particles were electrically charged,
then particles more distant from each other could interact. Including these interactions
(forces) would lead to non-zero terms farther from the matrix diagonal.
There are many techniques for solving Eq. (5.4), including those we developed in Chap-
ter 4. However, the methods developed in Chapter 4 become increasingly inefficient as the
number of particles N grows and the sizes of the matrices grow with them. For greatest
efficiency, we need to take advantage of some special properties of the matrices m and k ,
including their symmetry. We defer that discussion for now, however, in favor of a simpler
method based on the physics of the problem. We shall return to the more general method
afterwards, however, as it allows us to solve a much broader range of problems.
5.1.1 Small N solutions
We can gain some insight into the problem of N point masses on a string by considering
the small N limit. Let’s start with N = 1. Because there is only one particle, there is but one
degree of freedom and one equation of motion, which we can get directly from Eq. (5.2):
2T
mÿ1 = − y1 , (5.5)
l
where we have set the positions of the ends of the string to zero: y0 = y2 = 0. This is the
√
equation for a simple harmonic oscillator with an oscillation frequency of ω(1)
α = 2T /ml.
Here the superscript refers to the number of particles on the string, in this case 1; the
subscript refers to the normal frequency, labeled in Greek alphabetical order. For this trivial
case, there is one degree of freedom and one normal mode.
Let’s consider the case of two particles on a string: N = 2. In this case, there are two
t
equations of motion, one each corresponding to particles s = 1 and s = 2 in Eq. (5.2),
T
mÿ1 = (−2y1 + y2 ) (5.6)
l
T
mÿ2 = (y1 − 2y2 ) , (5.7)
l
where y0 = y3 = 0. Based on our experience with normal modes in Chapter 4, it is pretty
easy to guess what the normal modes must look like: a slow (α) mode where both particles
oscillate together (in phase) and a fast (β) mode where the two particles oscillate in opposi-
tion (exactly out of phase). These two modes are depicted in Fig. 5.2. Thus we would guess
that the normalized normal coordinates are given by qα = √12 (y1 + y2 ) and qβ = √12 (y1 − y2 ).
Exercise 5.1.1 By adding and subtracting Eqs. (5.6) and (5.7), show that you can obtain
two decoupled equations of motion where the normal coordinates are given by qα =
√1 (y1 + y2 ) and qβ = √1 (y1 − y2 ). What are normal frequencies of the two normal
2 2
modes?
For the case of three particles on a string, N = 3, the equations of motion become
T
mÿ1 = (−2y1 + y2 ) (5.8)
l
T
mÿ2 = (y1 − 2y2 + y3 ) (5.9)
l
T
mÿ3 = (y2 − 2y3 ) , (5.10)
l
where y0 = y4 = 0. We already treated this problem in §4.2.4, but even if you do not recall
N =1 n=1 N =3 n=1
N =2 n=1 N =3 n=2
N =2 n=2 N =3 n=3
t
Fig. 5.2 Normal modes for N masses on a string for N = 1, 2, 3.
t
the exact results, you can probably guess the qualitative form of the particle oscillations
corresponding to the three normal modes. You can compare your intuition against the Fig.
5.2, which shows diagrams of the various normal modes for the cases where there are 1 to
3 particles.
Exercise 5.1.2 Write down the equations of motion for the case where there are four
equally spaced particles of mass m on a string. By extending the diagrams in Fig.
5.2, try to guess the normal coordinates for the case where there are four particles.
Test your guesses by forming sums of from the linear combinations of y1 , y2 , y3 , and
y4 that form the normal coordinates. A correct guess should lead to an equation of
motion that is a function of a single normal coordinate.
5.1.2 Extrapolation to large N solutions
Returning to the general problem of N particles on a string, we begin by assuming a trial

solution to Eq. (5.2) of the usual form:
y s = a s e−iωt . (5.11)
Here y s and a s represent the instantaneous transverse displacement and the amplitude,
respectively, of the sth particle. By adopting this trial solution we are assuming, as we
always do in normal mode analysis, that all particles are oscillating at the same frequency
ω for a given normal mode. Substituting Eq. (5.11) into Eq. (5.2) yields
2T
!
T 2 T
− a s−1 + − ω m a s − a s+1 = 0 (5.12)
l l l
√
Dividing through by m and defining the characteristic frequency ω0 = T /ml, this be-
comes
−ω20 a s−1 + 2ω20 − ω2 a s − ω20 a s+1 = 0 .

(5.13)
Equation (5.13) actually represents N equations, one for each value of s, subject to the
boundary conditions that both ends are fixed at zero displacement: a0 = aN+1 = 0. While
solving this set of N coupled equations is straightforward for small values of N, it becomes
quite difficult for large values of N. Therefore, to assist us in our search for a solution, we
take a clue from the low-N solutions discussed above, which suggest a spatially oscillating
solution where the wavelength of the oscillation depends on the mode. So we start by
assuming a trial solution,1 say of the form
a s = Ae−iφ eiκs , (5.14)
where we take A to be real and explicitly include the possibility of non-zero phase φ. The
parameter κ is unknown at this point but its value clearly determines the wavelength of the
oscillation.
1 Physicists often call such a trial solution, which is generally an informed guess, an ansatz, which is a German
word meaning approach or attempt.
t
Let’s see if our trial solution Eq. (5.14) is consistent with Eq. (5.13). Substituting a s =
Aei(κs−φ) into Eq. (5.13) yields
−ω20 e−iκ + 2ω20 − ω2 − ω20 eiκ = 0 ,

(5.15)
where we have canceled out the common factors of Aei(κs−φ) . Solving for the eigenvalue ω2
gives
ω2 = 2ω20 − ω20 eiκ + e−iκ

(5.16)
= 2ω20 (1 − cos κ) (5.17)
κ
= 4ω20 sin2 , (5.18)
2
where we have used the identity 1 − cos κ = 2 sin2 (κ/2). Thus our trial solution Eq. (5.14)
is consistent with Eq. (5.13), with the eigenvalues given by Eq. (5.18). Equation (5.18)
applies for all boundary conditions, although, as we shall see, the set of values allowed
for κ is determined by the boundary conditions. At small values of κ, we can approximate
sin(κ/2) ' κ/2, so that Eq. (5.18) becomes
ω ' ω0 κ , (5.19)
√
where ω = T /ml: for small κ, the frequency ω increases linearly with κ.
Now let’s apply the boundary conditions at the left and right ends of the string. We start
with the simplest boundary conditions, namely that the ends of the string are fixed so that
a0 = aN+1 = 0. We first apply the boundary condition for s = 0, explicitly noting here that
we retain only the real part of our complex trial solution Eq. (5.14) This gives
a0 = Re(Ae−iφ ) = A Re(cos φ − i sin φ) = 0 , (5.20)
which is satisfied for φ = π/2. With φ = π/2, our trial solution Eq. (5.14) becomes
a s = Re Aei(κs−π/2) = A Re eiκs e−iπ/2

(5.21)
= A Re[(cos κs + i sin κs)(−i)] (5.22)
= A sin κs , (5.23)
which clearly satisfies the boundary condition at the left end of the string that a0 = 0. We
now apply the boundary condition at the right end of the string, which from Eq. (5.23)
gives
aN+1 = A sin κ(N + 1) = 0 . (5.24)
Equation (5.24) is satisfied only if κ(N + 1) = nπ, where n = 1, 2, 3, .... This gives
nπ
κn = where n = 1, 2, 3, . . . , (5.25)
N+1
where we have added a subscript n to indicate that each normal mode is characterized by
its own κn . Substituting the expression for κ into Eq. (5.18), we obtain an expression for
the normal mode frequencies
nπ
ωn = 2ω0 sin , (5.26)
2(N + 1)
t
ωn
2ω0
ω0
n
0 10 20 30 40 50 60 70
t
Fig. 5.3 Normal frequencies vs. mode number n for N = 16 identical masses on a string. Dark
circles denote modes for 1 ≤ n ≤ 16.
where we have added a subscript to ωn to denote that n = 1, 2, 3, . . .

While Eq. (5.24) is satisfied for any integer n, there are only N different normal modes
and normal frequencies. Figure 5.3 shows a plot of the (positive) normal frequencies as
a function of mode number n for n ≥ 0 as given by Eq. (5.26). There are N distinct
frequencies between n = 1 and n = N. There are no new or values of ωn outside this range;
in fact, as you can see from Fig. 5.3 or show algebraically from Eq. (5.26), ωn+2s(N+1) =
ω2s(N+1)−n = ωn for 0 ≤ n ≤ N + 1 where s is any integer. As we show below, none of the
particles move for the modes n = 0 and n = N+1. Thus, although mathematically n can take
on any integral value, there are only N distinct normal modes with finite displacements,
those corresponding to 1 ≤ n ≤ N.
The amplitude dependence given by Eq. (5.14) is for a single normal mode. Because
there are N different normal modes indexed by n, we add a subscript n to our equation for
the amplitude of the sth particle for the nth normal mode:
nπs
ans = An sin , where n = 1, 2, 3, ..., N . (5.27)
N+1
For a given normal mode n, Eq. (5.23) gives the amplitudes a s of the different particles
indexed by s = 1, 2, . . . , N, which are just the (unnormalized) entries in the corresponding
eigenvector. Thus, the eigenvector for the nth normal mode is
 sin κn 

 
2  sin 2κn 

r  
nπ
ẽn =  .  , κn = , where n = 1, 2, 3, ..., N, (5.28)
N + 1  ..  N+1
sin Nκn
 
√
with the prefactor 2/(N + 1) providing the normalization for the eigenvector.
The full general solution for the motion of the N masses is just a sum over all N eigen-
t
modes
N
X
y(t) = ẽn An e−iωn t , (5.29)
n=1
which can be expressed equivalently as

N
X
y(t) = ẽn (Bn sin ωn t + Cn cos ωn t) , (5.30)
n=1
where ωn is given by Eq. (5.26) and ẽn is given by Eq. (5.28). The instantaneous displace-
ment y s (t) of an individual particle is the sth component of y(t) is given by y s (t) = eTs · y(t),
where e s is a vector with an entry of 1 at the sth position and zeros everywhere else. Thus
N
X
y s (t) = eTs · ẽn (Bn sin ωn t + Cn cos ωn t) (5.31)
n=1
r
2
= sin κn s (Bn sin ωn t + Cn cos ωn t) (5.32)
N+1
where we have used Eq. (5.28) for ẽn and κn = nπ/(N + 1).
Exercise 5.1.3 Verify that the normalization of the eigenvectors given in Eq. (5.28)
does indeed give eigenvectors of length 1.
Normal modes 1 to N and beyond

It is instructive to examine the normal modes for at least one value of N. Let’s look at
N = 5, that is, the case of five particles of mass m evenly spaced on a string. The instan-
taneous displacements for each mode are shown for modes n = 1 to 12 in Fig. 5.4 (with
Bn = 0 and Cn = 1). First we note that modes 6 and 12 are null modes for which there are
no displacements. These represent trivial solutions (no motion) and thus are ignored. Next,
notice that modes 7–12 are the same as modes 1-5, but in reverse order and with a trivial
180◦ phase difference. You can also check that modes with the same pattern of displace-
ments oscillate at the same frequency, for example, ω5 = ω7 ' 1.93ω0 . Thus we see that
for N = 5, there are only 5 independent modes, which means we only need to consider
mode numbers n between 1 and 5. This result is quite general: although there are solutions
for n > N, they merely repeat the modes between 1 and N, so we do not include them in
our solutions. The eigenvectors for N = 5, obtained from Eq. (5.28), are
 1   1   1   1   1 
 2 √3   2   √3   2   2 √3 
 1   1 
 1   0  − 1 
 
 2         − 2 
2 2
 √1   √1   1 
ẽ1 =  3  , ẽ2 =  0  , ẽ3 =  3  , ẽ4 =  0  , ẽ5 =  √3  . (5.33)
   
 1   1   1   1 
 0 
 
 2  − 2   2   − 2 
 √1   1  √1   1  √1 
−2 −2
2 3 3 2 3
We also note here that small values of κ, which occur at small values of n are associated
with the longest wavelength normal modes. This is quite generally true.
t
n=1 n=7
n=2 n=8
n=3 n=9
n=4 n = 10
n=5 n = 11
n=6 n = 12
t
Fig. 5.4 Normal modes for five particles on a string (N = 5) shown when each particle is at its
maximum amplitude (dark) and at different times during each cycle (light). Modes 7-11
are not new modes but merely repeat modes 1-5, but in reverse order. Modes 0 (not
shown), 6, and 12 are null modes that have zero amplitude (i.e. they do not oscillate).
Exercise 5.1.4 What are the eigenvectors for the case N = 4? Once you find them,
explicitly show that they are orthogonal to each other.
t
An example with initial conditions
Equation (5.30), with the frequencies given by Eq. (5.26), represents the general solution to
the eigenvalue problem. To find the special solution for a specific case, we need to specify
the initial conditions, meaning the initial positions and velocities of all N particles. With
this information, the amplitudes {Bn } and {Cn } appearing in Eq. (5.30) can be determined
for any system of N particles on a string. Once the {Bn } and {Cn } are known, Eq. (5.30)
provides a complete description of the subsequent motion of the particles.
Suppose, for example, we give an sudden positive impulse to the second mass and a
sudden negative impulse to the fourth mass in a string of five identical masses. The initial
conditions are then y s (0) = 0 for all s and ẏ2 (0) = −ẏ4 (0) = 30 while ẏ s = 0 for s = 1, 3, 5.
The positions at t = 0 are all zero. From Eq. (5.30),
N
X
y(0) = Cn ẽn = 0 , (5.34)
n=1
which means that Cn = 0 for all n, which should not surprise you as this simply means that
the cosine term cannot contribute to the solution since all the particles have zero displace-
ment at t = 0.
Similarly, applying the initial conditions on the velocities by taking the time derivative
of Eq. (5.30), we have
 0 
 
 1 
 
XN
ẏ(0) = Bn ẽn = 30  0  (5.35)
 
n=1
 
−1
0
 
Each of the coefficients Bn can be isolated and determined by taking the dot products with
the five different eigenvectors of Eq. (5.33):
 1 T
 2 √3   0 
 
 1 
 1 
 
N N  2 
1 1
!
T
X
T
X  √1 
ẽ1 · ẏ(0) = ωn Bn ẽ1 · ẽn = ωn Bn δ1n = ω1 B1 =  3  · 30  0  = 30 − =0
 
n=1 n=1
 1    2 2
 2  −1
 √1  0
 
2 3
⇒ B1 = 0
 1 T
 0 
 
 2 
 1 
 1 
 
N N
1 1
X X  2  !
ẽT2 · ẏ(0) = ωn Bn ẽT2 · ẽn = ωn Bn δ2n = ω2 B2 =  0  · 30  0  = 30 + = 30
   
n=1 n=1
 1    2 2
− 2  −1
− 12 0
 
⇒ B2 = 30 /ω2
Continuing this procedure by similarly calculating ẽTn · ẏ(0) for n = 3, 4, 5 yields B1 = B3 =

t
B5 = 0 and B4 = −30 /ω4 . Combining this with the previous result that Cn = 0 for all n, the
complete solution follows from Eq. (5.30)
 1   1  
 2   2  
 1  − 1  
1 1 30  2  1  2 
" # 
y(t) = 30 ẽ2 sin ω2 t − ẽ4 sin ω4 t =  0  sin ω2 t − √  0  sin ω4 t
ω2 ω4 ω0  1  3  1  
− 2   2  
− 12 − 21
√
 sin ω0 t − √13 sin( 3ω0 t) 
 
y1 (t)
 
1
√
 sin ω0 t + √3 sin( 3ω0 t) 
 
y (t) 
 2  3 
y3 (t) = 0  0  , (5.36)

 2ω0  √
y4 (t) − sin ω0 t − √13 sin( 3ω0 t)
 
√
y5 (t) − sin ω0 t + √13 sin( 3ω0 t)
   
√
where from Eq. (5.26) ω2 = ω0 and ω4 = 3ω0 .
Exercise 5.1.5 Verify by explicit calculation that the initial conditions for the positions
and velocities of all five particles are satisfied by the solution given by Eq. (5.36).
The solution verifies some things we could have predicted from the beginning. Because
the initial conditions were antisymmetric about the central (s = 3) mass, only normal
modes having this same symmetry can contribute to the solution: only modes n = 2 and n =
4 are antisymmetric with respect to the central mass, so only they can be non-zero. Because
the displacements are all initially zero, Cn = 0 for all n. These simple considerations tell us
right from the beginning that only B2 and B4 can be non-zero. So we could have proceeded
directly to calculating their values, setting all the other coefficients to zero, based on these
simple ideas employing the symmetries of the initial conditions. Do not hesitate to employ
these kinds of ideas. It’s important to remember that shrewd thinking can often replace
tedious (and error-prone) calculations.
Indexing particles by distance

The horizontal physical distance along the string of a particle with index s is x s = sl. Using
this fact, we can replace s in in Eq. (5.30) by x s /l and rewrite our solution in terms of the
particle coordinate
N N
nπx s
X X !
y s (x s , t) = yns (t) = sin (Bn sin ωn t + Cn cos ωn t) (5.37)
n=1 n=1
(N + 1)l
N
X
= (sin kn x s )(Bn sin ωn t + Cn cos ωn t) , (5.38)
n=1
where
nπ
kn = . (5.39)
(N + 1)l
t
We recognize kn as the wave vector of the nth mode, with a corresponding wavelength of
2(N + 1)l/n.
We can also express the normal frequencies from Eq. (5.26) in terms of kn
1
ωn = 2ω0 sin kn l . (5.40)
2
For the lowest n modes where n N, we can make the approximation sin 21 kn l ≈ 12 kn l, in
which case Eq. (5.40) becomes
ωn ' ω0 kn l . (5.41)
Thus, for small n, which corresponds to the longest wavelength normal modes, the fre-
quency increases linearly with the wave vector kn . In §??, we show that this linear rela-
tionship between ωn and kn is consistent with a wave moving with a velocity vn = ωn /kn =
√
ω0 l = T l/m. As we shall see, the fact that frequency ceases to increase linearly with kn
(and n) has important consequences for the propagation of pulses in such a system, but we
defer discussion of wave velocity and pulse propagation to §??.
5.2 Identical masses coupled by springs
The problem we have been considering thus far, N masses on a string under tension, is
formally similar to the problem of N masses coupled by springs, illustrated in Fig. 5.5.
The spring problem does introduce some interesting new features, however, so we discuss
it briefly.
5.2.1 Masses coupled by springs with fixed ends
We first consider a set of N masses coupled by N + 1 springs where the ends of the first
and last springs are fixed. The equations of motion are given by
m ẍ1 = k(−2x1 + x2 ) (5.42a)

m ẍ2 = k(x1 − 2x2 + x3 ) (5.42b)
..
.
m ẍN−1 = k(xN−2 − 2xN−1 + xN ) (5.42c)
m ẍN = k(xN−1 − 2xN ) . (5.42d)
This set of equations has the same form as those in Eq. (5.3) and has exactly the same set
of solutions provided we replace coordinates yi with xi . That is, if we proceed with a trail
solution of the form
x s = a s e−iωt , (5.43)
141 Identical masses coupled by springs
t
x1 x2 xN
k k k k k k k
m m m m m m
t
Fig. 5.5 N spring-coupled masses with both ends fixed.
similar to Eq. (5.11), we obtain exactly the same equation for the amplitudes a s given
√
by Eq. (5.13), provided we replace the characteristic frequency ω0 = T /ml with ω0 =
√
k/m. Therefore, the normal modes have exactly the same form as we encountered for the
masses on a string under tension.
The most physically significant difference between the two problems is that here the
oscillations are longitudinal, or along the line connecting the masses, whereas of for the
masses on a string under tension, the oscillation are transverse, or perpendicular to the line
connecting the masses.
5.2.2 Periodic boundary conditions
Thus far we have studied the normal modes of a system of N masses connected by strings
or springs for one set of boundary conditions: fixed ends. It is instructive to examine what
happens for other boundary conditions. In Problems xx and yy at the end of the chapter
examine, you can examine what happens when there is no spring tethering the mass to a
wall at one or both ends of the line of spring-coupled particles.
Here we consider the very interesting case of periodic boundary conditions, that is,
boundary conditions where particle s + N is constrained to have exactly the same displace-
ment as particle s:
x s = x s+N . (5.44)
This boundary condition is equivalent to wrapping the line of springs around in a ring so
that the sth and the (s + N)th particles are actually identical, as illustrated in Fig. 5.6 for the
case where N = 6.
The normal mode analysis of this problem begins as usual by assuming temporally os-
cillating solutions of the form of Eq. (5.43). This leads to same amplitude equation given
by Eq. (5.13). Once again, guided by the expectation that the solutions will be spatially
m k m
xN x1
k k
m x2 m
k k
m k m
t
Fig. 5.6 N spring-coupled masses with periodic boundary conditions.
t
oscillatory, we assume solutions of the form
a s = Aeiκs , (5.45)
where in this case we allow A to be complex and thus to include the phase information.
Applying the periodic boundary condition of Eq. (5.44) gives
Aeiκs = Aeiκ(s+N) (5.46)
Multiplying through by e−iκs and canceling the common factor of A gives
eiκN = 1 , (5.47)
which is satisfied only if κN = 2πn, where n is an integer. This means that
2πn
κn = . (5.48)
N
where we expect that there are N independent normal modes, as there are N particles, each
with one degree of freedom.
The frequencies of the normal modes are still given by Eq. (5.18), but with κn given by
Eq. (5.48). Thus
πn
ωn = 2ω0 sin , (5.49)
N
with n = 0, 1, 2, . . . , N − 1. We could also choose n = 1, 2, 3, . . . , N, as modes 0 and N are
actually the same mode, as they must be owing to the periodic boundary conditions.
The mode corresponding to n = 0, or equivalently n = N, is particularly interesting.
From Eq. (5.49), we see that ωn=0 = ωn=N = 0. This is a zero-frequency mode in which all
the particles simply translate in the same direction (or rotate in the same direction if you
picture the system as a ring of particles as in Fig. 5.6). You can also verify from Eq. (5.45)
that a s = A; all particles have the same amplitude and move in the same direction.
The remaining normal modes are pretty much what you might expect in that they corre-
spond to oscillating solutions with finite frequencies and wavelengths, but with one curious
property. For the case we treated previously of fixed ends, κ ran from 0 to π (see Eq. (5.25)),
but for this case with periodic boundary conditions, κ runs from 0 to 2π (see Eq. (5.48)).
What is going on? It turns out that the eigenmodes for periodic boundary conditions are
traveling waves that can move either to the right or the left, whereas the eigenmodes for
systems with two ends are standing waves. With traveling waves, the whole waveform
moves to the left or the right (we will show this below), whereas for standing waves, the
waveform simply oscillates at its eigenfrequency without moving to the left or the right.
Physically, standing waves arise when there are boundaries that can reflect waves. For pe-
riodic boundary conditions, there are no boundaries and thus no reflections, so traveling
waves become possible.
To gain deeper understanding, we first note that from Eqs. (5.43) and (5.45), our eigen-
functions are of the form
x s = An e−iωn t eiκn s . (5.50)
In particular, x s ∼ eiκn s = ei(κn +2π)s , since ei2πs = 1 for every (integer) s. So rather than
143 Identical masses coupled by springs
t
ωn
2ω0
ω0
κn
−π 0 π 2π
t
Fig. 5.7 Modes of N = 16 spring-coupled masses with periodic boundary conditions: Small
dark squares show modes for periodic boundary conditions with 0 ≤ κn < 2π; Larger
gray squares show modes for periodic boundary conditions with −π ≤ κn < π; Gray
circles show modes for fixed-end boundary conditions with 0 ≤ κn < π. Note that there
is a mode at κn = 0 for periodic boundary conditions only.
restricting κ s to lie between 0 and 2π, we could equally well restrict it to lie between −π
and π, or equivalently, restrict N to lie between −N/2 and N/2. Figure 5.7 compares plots
of ωn vs. κn for periodic boundary conditions with 0 ≤ κn < 2πxth −π ≤ κn < π. It also
shows ωn vs. κn for boundary conditions with fixed ends treated in §5.1.2.
Rearranging terms, we can rewrite Eq. (5.50) as
x s (t) = An ei(κn s−ωn t) (5.51a)
= |An | cos(κn s − ωn t − δn ) , (5.51b)
where δn is the phase of the complex An . This function has a maximum whenever
κn s − ωn t − δn = 0 (5.52)
or, solving for s, for s = s p where
ωn
!
s p = s0n + t, (5.53)
κn
where s0n = δn /κn . Equation (5.53) says that the peak position s p moves to the right linearly
in time for κn > 0 and to the left for κn < 0 (since ωn ≥ 0). Thus we see that Eq. (5.51)
describes a sinusoidal wave moving to the left or to the right depending on the sign of κn .
It is instructive to return to the fixed boundary conditions. In that case the time dependent
position of the masses connected by springs is the same as that as for the masses connected
by strings given by Eq. (5.32). For notational simplicity, we absorb the normalization factor
into the (real) amplitude An and rewrite the temporal part in terms of cosine function and
a phase δn , while retaining the spatial dependent part sin κn s, with κn = nπ/(N + 1), as we
t
must, to satisfy the boundary conditions. Thus, the full solution is
x s (t) = An sin(κn s) cos(ωn t − δn ) (5.54)

1
= An [sin(κn s − ωn t + δn ) + sin(κn s + ωn t − δn )] , (5.55)
2
where we have used some trigonometric product identities. From Eq. (5.55), we see that
the solution for fixed boundary conditions consists of two sine waves, one, sin(κn s − ωn t +
δn ), propagating to the right and another, sin(κn s + ωn t − δn ), propagating to the left. As
noted above, these counter-propagating waves arise from the waves being reflected at the
boundaries.
The periodic boundary conditions we introduced above are often used in modeling phys-
ical systems, particularly when we want to focus on the behavior of the bulk material and to
avoid the effects of boundaries. Here we see that using periodic boundary conditions allows
us to explore propagating modes without concerning ourselves with the effects of reflec-
tions. We will find periodic boundary conditions useful when we study the propagation of
pulses.
5.3 Alternating masses
In our study of vibrations of a small number of coupled oscillators, we have considered

situations in which different masses are coupled by strings or by springs of different spring
constants. As the number of particles grows, the number of arrangements of strings or
springs and masses also grows, and it becomes impossible to obtain general expressions
for the normal modes that apply to all systems. There are, however, special arrangements
of different strings or springs and masses that are physically important. Solids composed
of more that one type of atom, for example, have periodically repeated unit cells composed
of two or more different atoms.
Here we consider one of the simplest models of such a system: an alternating sequence
of two unequal masses M and m evenly distributed along a string under tension T , as
shown in Fig. 5.8. We take the distance between unit cells—the repeat distance—to be l.
Therefore, for evenly distributed particles, the distance between masses M and m is l/2.
We label the cell number by s and denote the transverse displacements of the two masses
M and m for each cell by η s and ξ s . Following the analysis of §5.1, the equations of motion
for the two masses are given by
d2 ηs η s − ξ s−1 ξ s − η s
" #
M 2 =T − + (5.56)
dt l/2 l/2
d2 ξs ξ s − η s η s+1 − ξ s
" #
m 2 =T − + , (5.57)
dt l/2 l/2
where, as usual, we consider the displacements to be small: η s , ξ s l. Collecting terms,
145 Alternating masses
t
t
Fig. 5.8 A sequence of alternating masses M and m on a string under tension. The repeat
distance of each unit cell is l while the distance between successive masses l/2
(vertical scale exaggerated).
these equations become
d2 η s 2T
M = ξ s−1 − 2η s + ξ s (5.58)

dt 2 l
d2 ξ s 2T
m 2 = η s − 2ξ s + η s+1 . (5.59)

dt l
Once again we assume the trial solutions that oscillate in space and time
1
η s = Aη ei[κ(s− 4 )−ωt] (5.60)
i[κ(s+ 14 )−ωt]
ξ s = Aξ e , (5.61)
where we take the position of the mass m to be at s + 41 and the position of mass m to be
at s − 14 , which reflects the l/2 separation between the masses.2 Substituting these into the
equations of motion above, canceling common the common factor of ei(κs−ωt) , and collect-
ing terms gives
2T
Mω2 Aη = 2Aη − (eiκ/2 + e−iκ/2 )Aξ
h i
(5.62)
l
2T
mω2 Aξ = −(eiκ/2 + e−iκ/2 )Aη + 2Aξ .
h i
(5.63)
l
Noting that eiκ/2 + e−iκ/2 = 2 cos(κ/2), these equations can be recast in matrix form
ω2mA = kA , (5.64)
where
M 0 2T 2 −2 cos(κ/2) Aη
! ! !
m= , k= , A= , (5.65)
0 m l −2 cos(κ/2) 2 Aξ
where k is a symmetric matrix, as expected. To obtain a symmetric dynamical matrix, we

2 There is a good deal of freedom in formulating this problem. You can leave out the factors of ± 14 entirely,
which will change the solutions by shifting the phases of the amplitudes Aη and Aξ . This also makes the
amplitude ratios (as well as the dynamical matrix) complex, but ultimately yields solutions physically
equivalent to the ones we obtain here.
t
perform the same maneuvers we used in §4.2.5, defining the matrix
√
M 0
!
1/2
m = √ , (5.66)
0 m
and then left multiplying Eq. (5.64) by m−1/2
ω2m −1/2m A = m −1/2k m −1/2m 1/2 A

(5.67)
ω2m 1/2 A = m −1/2km −1/2 m 1/2 A ,

(5.68)
where we have inserted the identity matrix I = m −1/2m 1/2 between parentheses on the right
hand side of Eq. (5.67). Equation (5.68) can be rewritten as
D A0 = ω2 A0 (5.69)
where
2ω2M −2ω M ωm cos(κ/2)
!
D =m −1/2
km =
−1/2
, (5.70)
−2ω M ωm cos(κ/2) 2ω2m
√
M Aη
!
A0 = m 1/2 A = √ , (5.71)
m Aξ
where
r r
2T 2T
ωM ≡ , ωm ≡ . (5.72)
Ml ml
The solution to the normal mode problem is obtained now by finding the eigenvalues and
eigenvectors of the dynamical matrix D .
5.3.1 Eigenfrequencies of alternating masses on a string
We determine the eigenvalues ω2 of the dynamical matrix D as usual by finding the roots
of the characteristic equation
det D − ω2I = (2ω2M − ω2 )(2ω2m − ω2 ) − 4ω2M ω2m cos2 (κ/2)
h i
(5.73)
κ
= ω4 − 2(ω2M + ω2m ) ω2 + 4ω2M ω2m sin2 = 0 . (5.74)
2
Solving the quadratic equation for ω2 , we obtain an expression for the eigenvalues
q
ω2 = ω2M + ω2m ± (ω2M + ω2m )2 − 4ω2M ω2m sin2 (κ/2) .

(5.75)
From our previous experience, we expect that the parameter κ appearing in Eq. (5.75) will
be allowed to take on only certain discrete values, depending on the number of repeating
units Nc and the boundary conditions. We defer that discussion to the next section, how-
ever, and examine the solutions, which consist of two branches ω± (κ). The lower frequency
branch ω− is known as the acoustic branch because in solids, which have vibrational modes
very similar to those explored here, the frequencies of this branch are typically near acous-
tic frequencies. The upper frequency branch ω+ is known as the optical branch because in
solids, these frequencies typically occur at optical (or infrared) frequencies.
t
ωn
optical
branch 4ω0
√
2ωm
3ω0 band gap
√
2ωM
2ω0
acoustic
branch
ω0
κn
−π −π/2 0 π/2 π
t
Fig. 5.9 The spectrum of normal frequencies for a repeating sequence of Nc pairs of unequal
masses (here M = 2m). The solid lines are plots of Eq. (5.75), while the squares show
the allowed frequencies for periodic boundary conditions for the case of Nc = 16. The
dashed line shows ω = ω0 κ.
The two solid lines in Fig. 5.9 show plots of the two branches ω± (κ) of Eq. (5.75). The
limits of these two solutions as κ → 0 and κ → π are informative:

ω0 κ
 as κ → 0
ω− →  (5.76)

√
 2 ω M as κ → π;

q
 2(ω2m + ω2M ) as κ → 0


ω2+ →  (5.77)

√
 2 ωm

 as κ → π .
where
s
ωm ω M T
ω0 ≡ q = . (5.78)
l (M + m)
2(ω2m + ω2M )
The solutions assume ωm > ω M , consistent with M > m; the two frequencies ω M and ωm
are interchanged if the opposite is true. Thus ω− (κ) tends linearly to zero with κ, yielding a
zero frequency modes at κ = 0. On
√ the other hand,√ω+ (κ = 0) is nonzero. At κ = π, there is
a gap in the spectrum: ω+ (π) = 2ωm > ω− (π) = 2ω M . The gap gets smaller as M → m,
and vanishes altogether when M = m. The physical origin of this gap will become clearer
in the next section where we examine the eigenfunctions.
5.3.2 Eigenvectors for periodic boundary conditions
To find the eigenvectors (normal modes) for Nc repeating cells (here with two particles per
cell), we need to apply boundary conditions. First we consider periodic boundary condi-
t
tions:
ηNc +s = η s (5.79)
ξNc +s = ξ s , (5.80)
where s = 1, 2, . . . , Nc . Applying the periodic boundary conditions as we did in §5.2.2,
1 1
ηNc +s = Aη eiκ(Nc +s− 2 ) = Aη eiκ(s− 2 ) ⇒ eiκNc = 1 (5.81)
iκ(Nc +s) iκs iκNc
ξNc +s = Aξ e = Aξ e ⇒e =1, (5.82)
which is satisfied only if κNc = 2πn, where n is an integer. As before,
2πn
κn = . (5.83)
Nc
In contrast to the previous case for which there was only one particle and thus one degree
of freedom and one normal mode per unit cell (repeat distance l), here there are two parti-
cles per repeat unit cell, and thus two degrees of freedom and two normal modes. This is
reflected in Eq. (5.75), which has two solutions—two frequencies—for each allowed value
of κ = κn , which are plotted as gray squares in Fig. 5.9. As in the previous case, we use our
freedom to choose the modes such that −π < κn ≤ π, or equivalently −Nc /2 < n ≤ Nc /2.
The two sets of masses M and m oscillate with an amplitude ratio Aξ /Aη , which can be
obtained from the equation
D − ω 2 I A0 = 0 .

(5.84)
Writing this out explicitly using Eqs. (5.70) and (5.71), this becomes
! √
2ω2M − ω2 −2ω M ωm cos(κ/2) M Aη
!
√ =0. (5.85)
−2ω M ωm cos(κ/2) 2ω2m − ω2 m Aξ
Evaluating the top row of Eq. (5.85) gives
Aη m  2ω M ωm cos(κ/2) 
r  
=  , (5.86)
Aξ M 2ω2M − ω2+


while using the bottom row gives

Aξ
r
M 2ω M ωm cos(κ/2)
!
= . (5.87)
Aη m 2ω2m − ω2−
These two equations give the same results for the amplitude ratios, as we have found in
previous normal mode problems, and can be evaluated using Eqs. (5.75) and (5.83) for
ω± and κn . However, it is convenient to use Eq. (5.86) for √
the optical branch because the
denominator in Eq. (5.87) diverges at κ = π where ω+ = 2ωm . Similarly, it is better to
use Eq. (5.87) for the acoustic branch. The expressions for these amplitude ratios do not
simplify when ωn± and κn are substituted into these equations and thus are best evaluated
using a computer. As usual, the absolute magnitudes of the eigenvectors are fixed by the
normalization condition.
The normalized eigenvectors, obtained from the trial solutions given by Eqs. (5.60) and
(5.61) together with the amplitudes in Eqs. (5.86) and (5.87), are plotted in Fig. 5.10. In
t
n=0 n=0
s s
4 8 12 16 4 8 12 16
n=1 n=1
s s
4 8 12 16 4 8 12 16
n=2 n=2
s s
4 8 12 16 4 8 12 16
n=3 n=3
s s
4 8 12 16 4 8 12 16
n=4 n=4
s s
4 8 12 16 4 8 12 16
n=5 n=5
s s
4 8 12 16 4 8 12 16
n=6 n=6
s s
4 8 12 16 4 8 12 16
n=7 n=7
s s
4 8 12 16 4 8 12 16
n=8 n=8
s s
4 8 12 16 4 8 12 16
t
Fig. 5.10 Eigenvectors for modes n = 0 through 8 for N = 16 particles with alternating masses:
M = 2, solid circles, M = 1; open circles. Periodic boundary conditions. Left column,
acoustic (ω− ) branch; right column, optical (ω+ ) branch. Modes −1 to −8 are the same
those from 1 to 8, but propagate to the left rather than to the right.
t
each mode, the wavelengths of the oscillations for each set of masses, M and m, are the
same. However, in the acoustic branch the oscillations of the two types of masses are in
phase while in the optical branch they are out of phase. This can be seen in Fig. 5.10
and deduced from Eqs. (5.86) and (5.87), which have negative and positive denominators,
respectively.
The low κ (or n) modes of the acoustic branch are sinusoidal. The smallest non-zero
κ (n = 1) acoustic mode has a wavelength equal to the periodicity; longer wavelength
disturbances cannot satisfy the periodic boundary conditions and thus are not allowed. Each
new acoustic mode inserts another wavelength into the Nc l periodicity distance, which is
dictated by the periodic boundary conditions, codified in Eq. (5.48).
As noted above, ω ' ω0 κ for small κ, which corresponds wavelengths much greater than
the unit cell width l. The characteristic frequency ω0 , which relates ω to κ, is the tension
T per repeat distance l divided by the mass M + m in the unit cell. This is exactly the same
relationship we obtained at small κ in Eq. (5.19) for the case where all the masses were
identical. Apparently, the system does not care how the mass is distributed within the unit
cell when κ is very small, that is, when the wavelength of the disturbance is much larger
than l.
A quick perusal of Fig. 5.10 reveals that wavelengths, that is the characteristic distances
associated with one oscillation in space, are shorter for the modes in the optical branch
compared to those in the acoustic branch. In general, these shorter wavelength oscillations
are associated with high frequencies. Note, however, that the wavelengths of the oscilla-
tions in the acoustic and optical branches of the κ = π (n = 8) mode are exactly 2l—two
unit cells—in both cases. Interestingly, in the acoustic mode, the light particles of mass
m do not move, as predicted by Eq. (5.87) for κ = π; only the heavier mass M particles
oscillate. We can estimate the frequency of the oscillation of any mass M particle in the
κ = π acoustic mode by noting that it looks like a single-particle oscillator on a string of
tension T with fixed endpoints a distance l/2 away. According the the analysis in §5.1.1,
the oscillation frequency is 2T /M divided by the length of the string on each side of the
mass, which here is l/2. Therefore, the oscillation frequency of the large masses in the
acoustic mode is
s r
2T /M 4T √
ω8− = = = 2 ωM , (5.88)
l/2 Ml
where the last equality follows from the definition of ω M in Eq. (5.72). For the optical
mode when κ = π, the heavy particles of mass M do not move; only the lighter mass m
particles oscillate. By the same analysis, the oscillation frequency of the small masses in
the optical mode is
r
4T √
ω8+ = = 2 ωm . (5.89)
ml
Equations (5.88) and (5.89) give exactly the same results that we obtained in Eqs. (5.76)
and (5.77) from the normal mode analysis in the limit that κ = π. The origin of the band gap
at κ = π, therefore, is just the fact that the system switches between two equal wavelength
modes where either only the heavier masses or only the lighter masses oscillate.
t
Band gaps, such as the one seen here in a periodic system of two different masses os-
cillating on a string, are quite generally observed in periodic oscillating systems where
there are two underlying microscopic oscillation frequencies. These can arise from dif-
ferent masses, as is the case here, or from different interaction strengths, which occurs for
example in a system of identical masses linked by springs with alternating spring constants,
which also exhibits a band gap. In that case, the band gap comes about, again for the mode
where κ = π, when the system switches between low frequency oscillations, in which only
the weaker springs are excited (compressed and expanded), and high frequency oscilla-
tions, in which only the stronger springs are excited. This case is explored in problem zz at
the end of the chapter.
More generally, band gaps are observed in virtually all periodic systems that have two
or more underlying microscopic modes of oscillation. Photonic band gaps are observed
in the propagation of light in periodic dielectric materials where the electromagnetic field
oscillates at lower frequencies in the regions of higher dielectric permittivity, and oscil-
lates at higher frequencies in the regions of lower dielectric permittivity. Electronic band
gaps are observed in semiconductors and insulators where the frequency (or energy) of
the electron wave function oscillates slower or faster in regions of either weak and strong
Coulomb force. Similarly, band gaps can be observed for virtually any kind of wave in any
periodic medium in which the periodic changes in the relevant material property changes
the frequency at which the wave oscillates.
Band gaps are always observed for oscillations in one (spatial) dimension. In two and
three dimensions, the occurrence of band gaps becomes less likely because the length of the
unit cell is different in different directions, which, in turn, changes the range of frequencies
over which band gaps in different directions occur. For a complete band gap to occur in a
two or three dimensional material, there must be some range of frequencies in common for
the band gaps in different directions.
Band gaps can also arise in non-periodic systems, are long as there are at least two un-
derlying microscopic oscillation frequencies. Because of the absence of periodicity, how-
ever, such systems become much more difficult to treat analytically. In general this means
that we must turn to numerical schemes for analyzing such systems and determining their
normal modes of oscillation.
5.3.3 Eigenvectors for fixed boundary conditions
Here we consider boundary conditions where the ends of the string are fixed. The particle
displacements η s and ξ s are shown in Fig. 5.11. Since the first particle on the left has a
mass M with coordinate η1 , the boundary condition that the left end be fixed is given by
ξ0 = 0 . (5.90)
Similarly, since the last particle on the right has a mass m with coordinate ξNc , the boundary
condition on the right is given by
ηNc +1 = 0 . (5.91)
t
ηs+1 ηN−1
ξ2 ηs
η2 ηs−1 ξN−1
ξ1 ξs ξs+1 ηN
η1 ξs−1 ξN
t
Fig. 5.11 Coordinates for alternating masses M and m on a massless string under tension T .
We use trial solutions

1
η s = Bη e−iφ ei[κ(s− 2 )−ωt] (5.92)
−iφ i[κs−ωt]
ξ s = Bξ e e , (5.93)
which are slightly different from the trial solutions Eqs. (5.60) and (5.61) we used previ-
ously in that we place the mass m particle with transverse displacement ξ s at s and the mass
M with transverse displacement η s at s − 21 , which maintains the proper distance between
the particles but makes application of the fixed boundary conditions simpler. We have also
written the complex amplitudes Aη = Bη e−iφ and Aξ = Bξ e−iφ in polar form with real am-
plitudes Bη and Bξ and an explicit phase φ. These choices do not change Eqs. (5.62) and
(5.63) or the expression for the eigenfrequencies Eq. (5.75). Moreover the amplitude ratios
give by Eqs. (5.86) and (5.87) are maintained:
Aη Bη e−iφ Bη
= = . (5.94)
Aξ Bξ e−iφ Bξ
Applying the boundary condition for cell s = 0, we have
ξ0 = Bξ e−iφ = 0 , (5.95)
which is satisfied if φ = π/2. With this choice of phase, Eq. (5.93) becomes
ξ s = Bξ sin(κs) , (5.96)
which ensures that η0 = 0. With this choice of φ
η s = Bη sin[κ(s − 21 )] (5.97)
Applying the boundary condition for cell s = Nc + 1:
ηNc +1 = Bη sin[κ(Nc + 1) − 12 ] = Bη sin[κ(Nc + 12 )] = 0 . (5.98)
This boundary condition is satisfied if
nπ
κ= 1
, (5.99)
Nc + 2
where n is an integer that runs from 1 to Nc .

Equations (5.96) and (5.97) together with Eq. (5.99) and the equations for the amplitude
t
n=1 n=1
s s
4 8 12 16 4 8 12 16
n=2 n=2
s s
4 8 12 16 4 8 12 16
n=3 n=3
s s
4 8 12 16 4 8 12 16
n=8 n=8
s s
4 8 12 16 4 8 12 16
n=9 n=9
s s
4 8 12 16 4 8 12 16
n = 15 n = 15
s s
4 8 12 16 4 8 12 16
n = 16 n = 16
s s
4 8 12 16 4 8 12 16
t
Fig. 5.12 Eigenvectors for modes n = 1-3, 8, 9, and 14-16 for N = 16 particles with alternating
masses: M = 2, solid circles, M = 1; open circles. Fixed boundary conditions. Left
column, acoustic (ω− ) branch; right column, optical (ω+ ) branch.
ratios and eigenfrequencies obtained in the previous sections, constitute a complete solu-
tion to the problem. Figure 5.12 shows the eigenvectors for seven of the sixteen normal
modes. There are several ways in which these modes contrast with the normal modes for
the case of periodic boundary conditions, the most obvious being that the string is fixed
with zero displacement at both ends. There is no κ = 0 mode, as a system with its end tied
down cannot translate. The smallest finite-κ mode, n = 1, has roughly twice the wavelength
of the n = 1 mode for periodic boundary conditions, as the system with fixed boundary
conditions is not constrained to repeat itself. With fixed boundary conditions, each new
t
mode inserts one half new wavelength at small κ instead of a new whole wavelength as for
periodic boundary conditions.
There is a band gap between the acoustic and optical bands for mode n = 16 with
κ = 3233 π ≈ κ/2, and the frequencies of the two modes are nearly the same as they are
for periodic boundary conditions. While the heavy particles oscillate with large amplitude
for the acoustic mode and the light particles oscillate with large amplitude for the optical
mode, just as for periodic boundary conditions, the amplitudes of the complementary light
and heavy particles in the respective acoustic and optical do oscillate with relatively small
but still finite amplitudes, in contract to the case of periodic boundary conditions where
these particles do not move.
Finally, and perhaps most importantly, the modes for fixed boundary conditions are
standing waves, neither moving to the left nor to the right.
5.4 Nonperiodic systems
In the previous sections we considered the normal modes of sequences of identical masses
and springs or of alternating masses (or springs). In this section we explore the normal
modes of systems that are not periodic. We consider two examples: a random sequence
of two different masses and a quasicrystalline sequence of two different masses. These
systems exhibit phenomena similar to periodic systems, including band gaps and small-κ
(long wavelength) sinusoidal modes. They also exhibit new phenomena, including highly
localized oscillations that are still a topic of active research.
5.4.1 Diagonalizing the dynamical matrix
Nonperiodic systems are difficult, and in many cases simply impossible, to treat analyti-
cally, and so we turn to numerical methods.
5.4.2 Random sequences of masses: Localization
In the previous section we considered the normal modes of an alternating sequence of

masses on a string. Here we explore the normal modes of a random sequence of masses
on a string. Oddly enough, the normal modes that emerge are not completely random but
exhibit certain generic features, which have important consequences for the propagation
of waves in a wide variety of systems. Our simple system provides a convenient point of
entry into the fascinating physics of such systems.
First let’s define what we mean by a random sequence of masses. By random, we mean
that we select without bias each mass along the string from a well-defined distribution of
masses. This distribution is defined at the outset of the selection process. For example, we
might choose masses from a uniform distribution of masses between some minimum mass
155 Nonperiodic systems
t
m1 and maximum mass m2 . Then the distribution function is:
1

if m1 < m < m2



P(m) =  m m1 (5.100)


2 −

0 otherwise .


The probability that a particular mass is between m and m + dm is P(m) dm. Integrating
over all m should give unity
Z m2
1
Z ∞
P(m) dm = dm = 1 , (5.101)
0 m1 m 2 − m1
as indeed it does.
Alternatively, we might choose masses from a Gaussian distribution of width σ about a
mean of m0 :
2 2
P(m) = P0 e−(m−m0 ) /2σ , (5.102)
where the constant P0 is fixed by the normalization condition that the integral over all m
must yield unity probability, as in Eq. (5.101). In either case you can imagine a process in
which masses are randomly selected out of a hat that contains a large number of particles
whose masses are distributed according to some well-defined distribution such as those
given above.
Here we consider a particularly simple distribution of masses: only two different masses,
m1 and m2 , but arranged in a random sequence. Problems at the end of the chapter explore
the normal modes of systems with other kinds of random distributions. For the present
case, we imagine that the probabilities of choosing a mass m1 or m2 are equal.
First we need to specify a random sequence for the two masses. We could do so by
flipping a coin: heads, choose m1 , tails, choose m2 ; repeat until you have a sequence of
32 masses. This is much too tedious, however, so we let a computer do it for us using an
algorithm known as a random number generator. We are not going to delve into how one
develops such an algorithm; you can read about that elsewhere.3 We simply assume that
we can find a random number generator—a computer function—suitable for our purposes.
Having generated a random sequences of masses, we need to solve the normal mode
problem. Once again, we turn to the numerical methods and the computer programs we
have developed. To solve the present problem, we need only to populate the diagonal ele-
ments of the mass matrix m with our random sequence of masses and then run the computer
program that diagonalizes the matrices and calculates the normal modes. We present the
results generated by such a numerical solution below.
Figure 5.13 shows the first two normal modes as well as normal modes 16 and 17 for
a random sequence of masses m1 and m2 where the total number of masses is 32. Here
we use the same physical parameters used in previous sections: m1 = 2.5 g, m2 = 10 g,
l = 4 cm, and T = 60 N. The lowest frequency modes look essentially the same as those
for the cases of a identical masses and alternating masses, respectively, as can be checked
by comparing Fig. 5.13 with Figs. (??) and (??). In the limit that the number of masses
3 Numerical Recipes by Press et al. provides a useful introduction to random number generators.
t
0.4
1
ys
0.0
8 16 24 32
0.4 particle number s
2
ys
0.0
8 16 24 32
−0.4 particle number s
0.4 16
ys
0.0
8 16 24 32
0.4 17
ys
0.0
8 16 24 32
t
Fig. 5.13 Normal modes for a random sequence of light particles (small circles) and heavy
particles (large circles). There are 32 total particles. Mode number is indicated in
upper right corner of each plot.
N 1, the oscillation frequencies of the low n modes for all three cases is given by
r
T l πn
ωn = , for n N , (5.103)
m̄ L
where m̄ is the average mass of a bead on the string and L is the length of the string.
Exercise 5.4.1 Starting from Eqs. (5.41) and (??), show that Eq. (5.103) holds for the
cases of identical masses and alternating masses on a string. Equation (5.103) also
holds the in the case of a prandom sequence of masses. Why? Note that Eq. (5.103)
can be rewritten as ωn = T /µ (πn/L) where µ = m̄/l is the mass per unit length of
the beaded string.
Modes 16 and 17 for our random sequences of masses are also shown in Fig. 5.13. In
contrast to the low-n low-freqency modes, these higher frequency modes do not resemble
the normal modes observed previously for the case of identical or for the case of alternating
masses. Moreover, modes 16 and 17 will look different from those depicted here for other
realizations of the system, that is, for other random sequences of the same masses m1 and
m2 . For example, Fig. 5.14 shows normal modes 16 and 17 for a different random sequence
of masses. This illustrates an important point: the particular shapes of the larger n normal
modes depends on the particular sequence of masses. By contrast, the low-n low-frequency
modes all look very similar and have similar normal frequencies.
Although the higher-n modes change significantly for different random sequences of
masses, they do exhibit some common features. In particular, a relatively small fraction
t
of the particles in modes 16 and 17 have large displacements, as shown in Figs. (5.13)
and (5.14). The region of large particle displacements is confined to a range of about 10-14
particles for modes 15 and 16. Such modes are said to be localized and are a generic feature
of all but the lowest frequency modes in systems with a random distribution of masses.
0.4 16
ys
0.0
8 16 24 32
0.4
17
0.0
ys
8 16 24 32
t
Fig. 5.14 Normal modes for a random sequence of light particles (small circles) and heavy
particles (large circles) that is different from that shown in Fig. 5.13. There are 32 total
particles. The mode number is indicated in the upper right corner of each plot.
32
0.4
ys
0.0
8 16 24 32
31
0.4
ys
0.0
8 16 24 32
30
0.4
ys
0.0
8 16 24 32
t
Fig. 5.15 High frequency normal modes for a the same random sequence of light particles and
heavy particles shown in Fig. 5.13. The mode number is indicated in the upper right
corner of each plot.
t
In general, as the mode number and frequency increases, the normal modes in a random
system become more and more localized. Figure 5.15 shows the three highest-n modes for
the original random sequence of masses used in Fig. 5.13 and we see that they are even
more localized than the n = 15 and 16 modes. We can characterize the spatial extent—
the localization—of a normal mode quantitatively by averaging over the square of the
amplitudes of the particles in a given mode. Because the energy is proportional to the
square of the particle amplitude, we are weighting according to where energy of the mode
is located. Thus, the average position of a normal mode is given by
N
X
s̄ = s y2s . (5.104)
s=1
Here we use normalized eigenvectors so that s y2s = 1 so no additional normalization is

P
needed. The spatial extent of a normal mode can then be characterized by
 N 1/2
X 2 2
` = l  (s − s̄) y s  .

(5.105)
s=1
The length ` gives the spatial extent of a normal mode, and is commonly known as the
localization length. In general, ` depends on the frequency of the mode and, to a lesser
extent, on the particular realization of the random sequence of masses. From the depiction
of normal modes in Fig. 5.13, we might expect that ` ∼ L for the low-n modes. By contrast,
for the From Figs. (5.14) and (5.15), we would expect ` L for the high-n high frequency
modes. Indeed, this trend is supported by Fig. 5.16(a), which shows how ` varies with
frequency for the system of 32 masses depicted in Figs. (5.13) and (5.15). As expected, `
decreases with ω at high frequencies but is relatively flat at low frequencies, with a value
of about one third of the systems size—about 10 particles out of 32. You can get a sense of
(a) 103 (b)

101
102
l
101
100
100
10−1 10−1
102 103 100 101 102 103
ω ω
t
Fig. 5.16 Localization length vs. frequency. (a) N = 32. Red data shows ` for the sequence of
32 masses used in Fig. 5.13; black data are for 20 different realizations of sequences
of 32 masses. (b) Data for N = 2048 masses. Gray line shows ω−2 for reference.
t
how much ` can vary from one random sequence to another from the smaller gray points
that show ` vs. ω for 20 different realizations of a random sequence of 32 masses. Clearly,
the fluctuations in ` from one realization to another can be significant. The large variation in
` from one realization to another is due in part to the relatively small number of particles in
the system, 32 in this case. Figure 5.16(b) shows data for a string with 2048 masses, where
the fluctuations in ` are smaller and the trend for ` to decrease with increasing frequency
is quite clear. In fact, the results for this larger system are good enough to suggest that
` ∼ ω−2 for the high frequency modes, as shown by the gray line in Fig. 5.16(b). The
localization length ` ceases to increase as ` approaches the system size L, just as it did for
the smaller system.
Now that we know something about the normal modes of our random system, let’s ex-
amine the normal frequencies. Figure 5.17 shows the normal frequencies as a function of
mode number for the system of 32 masses (dark circles). Interestingly, Fig. 5.17 reveals
that there appear to be a number band gaps, although the gaps are smaller than the single
gap we found for the case of alternating masses. Figure 5.17 also shows the normal fre-
quencies for a different random sequence of masses (light circles). The normal frequencies
of this second sequence of are very nearly the same for the low-frequency modes but show
significant differences at the higher frequencies. The two realizations of the system have
band gaps at similar, but not exactly identical frequencies. This is not surprising as the two
sequences are different.
Exercise 5.4.2 Why are the normal modes of the two different realizations of a random
sequence of masses very nearly the same at low frequencies but fairly different at
other frequencies, particularly those in the mid-range? [Hint: See Exercise 5.4.1]
1600
1400
1200
1000
ω (rad/s)
800
600
400
200
0
8 16 24 32
mode n
t
Fig. 5.17 Normal modes for a random sequence of 32 spheres: dark circles, random sequence
of masses used in Fig. 5.13; light circles, random sequence of masses used in Fig.
5.14.
t
Examining Fig. 5.17 reveals that for both realizations of the system, there are no normal
modes in the frequency ranges from about 920 s−1 to 1130 s−1 and from about 1200 s−1 to
1350 s−1 . Is this a generic feature of all realizations of this system or just a coincidence for
these two realizations? To address this question, we need to examine the normal modes for
a large number of random systems that are statistically equivalent. By statistically equiva-
lent, we mean that for each realization, the sequence of masses is chosen randomly from
the same probability distribution. Such a collection of systems is called an ensemble, and
averaging some physical quantity of interest over this ensemble is called an ensemble av-
erage.
To answer the question about band gaps posed above, we calculate the normal frequen-
cies for a large number of realizations of this system and make a histogram of the fraction
of the total number normal frequencies between ω and ω+δω, where δω is fixed to be some
convenient small frequency interval. We call frequency-dependent height of the histogram
the density of normal modes or, more commonly, the density of states D(ω). Figure 5.18
shows the density of states calculated for an ensemble of 32 particles and for an ensemble
of 2048 particles. The density of states is essentially the same for the two systems, with
D(ω) being more finely grained due to the larger number of modes. In both cases, there is
a remnant of the band gap just below ω1 , the lowest frequency mode of the optical band
(a) ω2 ω1 ωt ωm
(ω )
0 200 400 600 800 1000 1200 1400 1600

ω
(b) ω2 ω1 ωt ωm
(ω )
0 200 400 600 800 1000 1200 1400 1600

ω
t
Fig. 5.18 Density of states D(ω) for a random sequence of masses m1 and m2 . (a) N = 32; (b)
N = 2048. Dark vertical lines indicate the frequencies ω1 , ω2 , and ωt corresponding to
the frequencies of the top of the acoustic branch and the bottom and top of the optical
branch, respectively, for an alternating sequence of masses m1 and m2 . The frequency
ωt corresponds to the maximum normal frequency consisting of all light masses m1 .
t
for a system of alternating spheres. Moreover, there is a surplus of states between ω1 and
ωt , the range of frequencies occupied by the acoustical branch for the case of alternating
masses. There are also states above ωt , the maximum frequency normal mode for the case
of alternating spheres. Problem ... at the end of the chapter explores the physical origins of
these features of the density of states.
5.4.3 A Fibonacci sequence of two masses: Quasiperiodicity
An interesting case that is intermediate between a periodic and a random sequence of

masses is a quasiperiodic sequence of masses. A simple example of a 1-dimensional qua-
sicrystal can be constructed from a Fibonacci sequence. You may be familiar with a Fi-
bonacci sequence of integers: 1, 2, 3, 5, 8, 13, 21, 34, 55, 89, . . .. The sequence starts with
0, 1 with the rule that the next element fi of the sequence is generated from the sum of the
previous two elements: fi = fi−1 + fi−2 .
We can use the same rule to create a binary sequence of two letters, say A and B. Starting
with the sequences f1 =A and f2 =AB, f3 = f2 + f1 =ABA, f4 = f3 + f2 =ABAAB, f5 =
f4 + f3 =ABAABABA, etc., where in this case the “+" operation denotes the concatenation
of two sequences. Among the many very interesting properties of the Fibonacci sequence
is that the pattern of letters A and B never repeats itself, somewhat like the digits of an
irrational number. However, the sequence of letters is not exactly random either, as we
shall see.
Here we consider a Fibonacci sequence of the same two masses m1 and m2 that we used
for the alternating sequence and for the random sequences of masses. This time, however,
we construct the sequence of masses the way we constructed the sequence of letters A and
B above: the letters A and B correspond to the masses m1 and m2 , respectively. Instead
of constructing our systems of 32 or 2048 particles, the number of particles will be a
Fibonacci number fn . This means that the number of masses m1 and m2 will be fn−1 and
fn−2 , respectively. Let’s start with N = 34, which means that there are 21 particles of mass
m1 and 13 particles of mass m2 arranged in the following sequence:
ABAABABAABAABABAABABAABAABABAABAAB. (5.106)
We solve the eigenvalue problem once again simply by substituting the sequence of
masses, in this case the Fibonacci sequence, into the mass matrix. Figure 5.19 shows the
resulting normal mode frequencies as a function of the mode number. Three obvious band
gaps are visible, just above modes 13, 21, and 26. Note that the number of modes between
each set of gaps in a Fibonacci number: 13, 8, and 5. Alternatively, we note that the three
gaps occur at 8, 13, and 21 modes, all Fibonacci numbers, counting right to left from the
highest frequency mode.
Figure 5.20 shows the two lowest frequency modes, which as in the cases for other
sequences of masses, are modes with wavelengths of 2L and L, where L is the length of
the string. Figure 5.20 also shows modes 13 and 14, the modes just below and just above
the lowest-frequency gap, respectively. In mode 13, the heavy masses m2 have the largest
amplitudes, while in mode 14, the light masses have the largest amplitudes with the large
t
1400
1200
1000
ω (rad/s)
800
600
400
200
0
8 16 24 32
mode n
t
Fig. 5.19 Frequencies of normal modes vs. mode number for a Fibonacci sequence of two
masses.
particles remaining close to their equilibrium positions. It is this sudden change in the mass
that oscillates without a corresponding change in the wavelength that leads to the band gap.
As the normal mode frequencies increase, the oscillations become more localized, al-
0.4
1
ys
0.0
8 16 24 32
0.4 particle number s
2
ys
0.0
8 16 24 32
0.4
13
ys
0.0
8 16 24 32
0.4
14
ys
0.0
8 16 24 32
t
Fig. 5.20 Normal modes for a Fibonacci sequence of two masses. The mode number is shown
in the upper right corner of each plot.
t
0.4
89
ys
0.0
16 32 48 64 80
0.4
88
ys
0.0
16 32 48 64 80
particle number s
−0.4
0.4
87
ys
0.0
16 32 48 64 80
t
Fig. 5.21 Normal modes for an 89-mass Fibonacci sequence of two masses. The mode number
is shown in the upper right corner of each plot.
though in a way that is quite distinct from the localized modes of the random system. As it
is difficult to appreciate the character of these high frequency modes for a system of only
34 particles, we illustrate the high frequency modes for a system of 89 particles. Figure
5.21 shows the normal modes for the three highest frequencies of a two-mass system of
89 particles arranged in a Fibonacci sequence. Mode 87 is similar to localized modes in a
random system in that nearly all the particles with large displacements are near each other.
Modes 88 and 89 are more typical, however, of high frequency modes in a Fibonacci se-
quence. In both cases there are two sets of well-separated localized oscillations. Thus, one
cannot say that those modes are localized in the same sense that high frequency modes for
a random distribution of masses are localized. Unlike the system of identical particles or
the system of alternating masses, however, the oscillations do not appear to extend across
the entire system. Thus, the character of the normal modes of the Fibonacci sequences,
with their well-defined band gaps but quasi-localized high frequency modes, seems to be
somewhat intermediate between an ordered system of alternating masses and a completely
disordered system of randomly distributed masses.
t
Problems
5.1 Consider a system of five particles of mass m evenly spaced on a massless string
fixed at both ends and stretched in a straight line to a tension T . The middle mass
is then pulled up so that the system has the initial configuration shown in the figure
below. The particles are initially all at rest.
y
0.06
0.04
0.02
0.00 x
1 2 3 4 5 6
− 0.02
(a) At time t = 0, the masses are released and begin to oscillate. Starting from Eq.
(5.30), find the subsequent motion of all five the particles. Hint: This problem
boils down to finding the ten coefficients {Bn } and {Cn } in Eq. (5.30). You can
set all but three of the ten coefficients to zero by considering the initial con-
ditions, including the symmetry of the initial configuration with respect to the
symmetries of the five normal modes indexed by n.
(b) Using your solution from part (a), plot the configurations of the particles at sev-
eral different times as the system oscillates. Does the system ever return to its
initial configuration? If so, when? If not, why?
5.2 Show that the limit that κ 1, Eq. (??) yields the following two solutions
s
ω1 ω2 T
ωa = κ + O(κ3 ) = κ + O(κ3 ) (5.107)
2(ω1 + ω22 )1/2
2 2(m1 + m2 )l
s
2T 1 1
!
ωo = (ω21 + ω22 )1/2 − O(κ2 ) = + − O(κ2 ) . (5.108)
l m1 m2
Therefore, as discussed in the text, the acoustic branch has a linear dispersion rela-
tion in the limit of small κ or long wavelengths. By contrast, the dispersion relation
for the optical branch is nearly independent of frequency and thus highly dispersive.
5.3 Starting from Eqs. (??) and (??) derive the following expressions for the normal
165 Problems
t
mode amplitude ratios
ω22 − ω2a
ηn /ξn = , acoustical branch (5.109)
ω22 cos(κ/2)
ω21 − ω2o
ξn /ηn = , optical branch , (5.110)
ω21 cos(κ/2)
and plot the results using the expressions for ωa and ωo given by Eq. (??). Set N = 32
for your plots and explain briefly in words what the plots signify for n = 1 and at the
band gap.
5.4 Show by writing y(x, t) in terms of its temporal Fourier transform Y(x, ω) (different
from Y(k, t)) that you can obtain a different equation
d2 Y(x, ω)
+ k2 (ω) Y(x, ω) = 0 , (5.111)
dx2
where in this case Y(x, ω) is the temporal Fourier transform of y(x, t).
5.5 Consider a linear string of masses all with the same mass m connected to each other
by springs whose spring constants alternate between k1 and k2 as shown in the fig-
ure below. The masses are constrained to move only in the horizontal direction (no
transverse motion).
(a) Write down the system of equations of motion for the masses and then express
them in the standard matrix form m ẍ = −k x. What are the mass and stiffness
matrices m and k ? Note: You will need to write down separate equations of
motion for for masses with even and odd values of s.
(b) Setting m = 10 g, k1 = 3 N/m, and k2 = 6 N/m, solve for the eigenvectors and
eigenvalues of the system for the case of N = 16 masses using the numerical
routines available on the NYUClasses course web site. Plot x s vs. s for modes
1-4, 7-10, and 13-16 (i.e. the eigenvectors) Plot all the normal frequencies as a
function of mode number from lowest to highest mode number. You should find
a band gap between modes 8 and 9. Compare your results to the case of two
different masses but identical springs discussed in the text. Explain physically
why there is a band gap by considering the modes just below and just above the
gap.
6 Fourier Analysis
The waves we encounter in nature, whether they are the ripples on the surface of a pool
of water or the sounds of a jazz ensemble, are seldom characterized by the simple single-
frequency and single-wavelength sine and cosine functions introduced in elementary ac-
counts of waves. The sound waves produced by even the simplest musical instrument have
a much more complex harmonic content, meaning they are made up of waves simultane-
ously oscillating a number of different frequencies. It is this rich harmonic content that
makes an A-note played on a violin sound different from the “same" A-note played on a
trumpet. Take a look at the waveforms shown in Fig. ??, which were taken from a violin and
a trumpet. The two waveforms have the same overall period, but look rather different, as
you might expect given the very different quality of the sounds the two instruments make.
The waveforms differ in their harmonic content, that is, in the spectrum of frequencies that
make up the sound we hear when each instrument is played. Fourier analysis attempts to
characterize mathematically, and rather precisely, just how the waveforms and harmonic
content of the trumpet and the violin differ.
The ripples on the surface of a pool of water present another kind of waveform, one that
is clearly a wave that oscillates up and down, but one that has a finite extent: it starts at
some some position in space, grows bigger, reaches a maximum, and then decays away.
This waveform is more-or-less sinusoidal, except its amplitude grows and decays in space.
If instead of taking a snapshot of the spatial undulations of the waveform, you monitored a
single point on the surface of the water, you would see it start to oscillate up and down in
time as the wave passed by. Once again, the amplitude of the waveform initially grows, this
time in time, reaches a maximum, and then decays away. How do we describe the harmonic
content this waveform? The waveform superficially resembles the beats we encountered
in our study of two coupled oscillators in Chapter 4. There we found that two coupled
oscillators could produce beats if both modes were simultaneously excited and if the two
modes had very similar frequencies. This provides a clue as to the harmonic content of
the waveform of the water surface wave: perhaps it consists of two or more sinusoidal
waves of nearly the same frequencies. This is in fact the case. The precise mathematical
description of such waves is once again the subject of Fourier analysis—in this case Fourier
transforms.
In this chapter, we study Fourier series and Fourier transforms. Fourier series can de-
scribe virtually any periodic function. However, Fourier series can only describe periodic
functions of infinite extent, that is, functions that repeat forever either in space or time. To
describe wave pulses, that is, waveforms of finite extent in space or time, we need a more
general approach, which involves Fourier transforms. In §6.1 we introduce Fourier series
and in §6.2 we introduce Fourier transforms.
166
167 Fourier series
t
Computational physics, which involves encoding information in digital form, puts fur-
ther demands on our mathematical description of waves. To address these demands, we
develop discrete Fourier transforms (DFTs) in §6.3 as well as algorithms for calculating
DFTs, including the famous fast Fourier transform (FFT) algorithms.
6.1 Fourier series
The basic idea of Fourier series is that virtually any periodic function can be represented
as an infinite series of sine and cosine functions. Fourier series are useful not only as
mathematical tools but also as a means for understanding the harmonic makeup of complex
non-sinusoidal periodic waveforms.
6.1.1 Sine and cosine Fourier series
In physics, we generally encounter functions that are either periodic in space or in time or
in both space and time. In this section, we explicitly deal with functions that are periodic in
space. However, we could equally equally well work in the temporal domain. Ultimately,
we will study functions that are periodic in both space and time.
Figure 6.1 shows two examples of periodic functions. Mathematically, we can express
the fact that a function f (x) is periodic in space by writing f (x + L) = f (x). We call L
the periodicity (or wavelength) of the function. Such a function can quite generally be
represented by an infinite series of sine and cosine waves:
∞ ∞
1 X X
f (x) = a0 + an cos kn x + bn sin kn x (6.1)
2 n=1 n=1
where kn = 2πn/L and n = 1, 2, 3, ... . For n = 1, k1 = 2π/L and the sine and cosine waves
have a wavelength of L. For n = 2, 3, ..., the sine and cosine waves have successively shorter
wavelengths L/n corresponding to the harmonics of the n = 1 sine and cosine waveforms,
as shown in Fig. 6.2.
For a given periodic function f (x) there is a unique set of coefficients an and bn that,
when used in Eq. (6.1), yield a Fourier series that converges to the desired periodic func-
tion. We can obtain the coefficients using a simple procedure. First, we multiply both sides
(a) 6 (b) 2
−8 −6 −4 −2 2 4 6 8
x −4 −3 −2 −1 1 2 3 4
x
−6 −2
t
Fig. 6.1 Periodic wavefunctions: (a) sawtooth wave with a period L = 4, (b) triangle wave with
a period L = 2.
168 Fourier Analysis
t
n=1 1 n=1 1
−L/2 L/2 −L/2 L/2

−1 −1
n=2 1 n=2 1
−L/2 L/2 −L/2 L/2

−1 −1
n=3 1 n=3 1
−L/2 L/2 −L/2 L/2

−1 −1
t
Fig. 6.2 Cosine (blue) and sine (red) waves for the first three terms n = 1 to 3 in a Fourier
series. The solid lines indicate the interval −L/2 ≤ x ≤ L/2.
of Eq. (6.1) by cos km x, where km = 2πm/L and m can be any positive integer, and then
integrate over the interval [−L/2, L/2]:
L/2 Z L/2
1
Z
f (x) cos km x dx = a0 cos km x dx
−L/2 2 −L/2
∞
X Z L/2
+ an cos kn x cos km x dx
n=1 −L/2
X∞ Z L/2
+ bn sin kn x cos km x dx . (6.2)
n=1 −L/2
The first integral on the right hand side of Eq. (6.2) is zero since cos km x has positive and
negative parts in equal measure on the interval (−L/2, L/2). The third integral on the right
hand side of Eq. (6.2) is also zero since the product sin kn x cos km x is an odd function in-
tegrated over an interval (−L/2, L/2) that is symmetric about x = 0. Finally, the second
integral is zero unless km = kn . If km = kn , then the integral is just 1/2, the average of
cos2 kn x, times L, the length of the interval of integration. Mathematically, we can summa-
rize the integrals over the sine and cosine functions by
L/2
L
Z
cos kn x cos km x dx = δnm (6.3)
−L/2 2
Z L/2
sin kn x cos km x dx = 0 (6.4)
−L/2
where δnm is the Kronicker delta, which is defined to be 1 if n = m and 0 if n , m.

Therefore, the only term in the sum over an on the right hand side of Eq. (6.2) that is
169 Fourier series
t
nonzero is the one where m = n. Thus, Eq. (6.2) becomes
Z L/2 ∞
X Z L/2
f (x) cos km x dx = an cos kn x cos km x dx
−L/2 n=1 −L/2
∞
X L
= an δnm
n=1
2
L
= am . (6.5)
2
We can solve this to obtain as expression for am . Once we have obtained the result, however,
the index label is arbitrary, so we can write the result using the index n:
L/2
2
Z
an = f (x) cos kn x dx . (6.6)
L −L/2
An equation for bn can be found using a similar approach by multiplying both sides of Eq.
(6.1) by sin km x and integrating over the interval [−L/2, L/2]. The result is
L/2
2
Z
bn = f (x) sin kn x dx , (6.7)
L −L/2
Finally, an expression for a0 can be obtained simply by integrating both sides of Eq. (6.1)
over the interval [−L/2, L/2]. The sine and cosine terms integrate to zero giving:
L/2
2
Z
a0 = f (x)dx . (6.8)
L −L/2
Before proceeding, we point out that we could have performed the integrals in Eqs. 6.6–
6.8 over any interval of length L and still have obtained the same results for the coefficients
an and bn . For example, we could have used the interval [0, L]. Try it out! Because any
waveform described by a Fourier series is periodic with a periodicity of L, the choice of
interval is up to you.
We are now in a position to use Eq. (6.1), together with Eqs. 6.6–6.8 for the Fourier
coefficients an and bn , to determine the Fourier series for a few example waveforms.
Let’s start by finding the Fourier series for the waveforms shown in Fig. 6.1. Finding the
Fourier series for a periodic waveform f (x) is, according to Eq. (6.1), simply a matter of
determining the coefficients an and bn , which we can do using Eqs. 6.6–6.8. To calculate the
integrals in Eq. (6.1), we need to specify the function f (x) only in the interval of integration
[−L/2, L/2], i.e. one period. For the waveform in Fig. 1(a), the interval of integration is
t
(a) 6 (b) 6
x x
−8 −6 −4 −2 2 4 6 8 −8 −6 −4 −2 2 4 6 8
−6 −6
(c) 6 (d) 6
x x
−8 −6 −4 −2 2 4 6 8 −8 −6 −4 −2 2 4 6 8
−6 −6
t
Fig. 6.3 Fourier series (black) for sawtooth function (red) including: (a) n = 1 term, (b) n ≤ 2
terms, (c) n ≤ 3 terms, and (d) n ≤ 4 terms.
[−2, 2] ( L = 4) where the function is f (x) = 3x. The coefficients are then given by
2 L/2 2 2
Z Z
a0 = f (x) dx = 3x dx = 0 (6.9)
L −L/2 4 −2
2 L/2 2πnx
Z !
an = f (x) cos dx (6.10)
L −L/2 L
1 2 πnx
Z
= 3x cos dx = 0 (6.11)
2 −2 2
Z L/2
2 2πnx
!
bn = f (x) sin dx (6.12)
L −L/2 L
1 2 πnx
Z
= 3x sin dx (6.13)
2 −2 2
−12 12
= 2 2 [πn cos(πn) − sin(πn)] = (−1)n+1 , (6.14)
πn πn
where we used the fact that cos(πn) = (−1)n and sin(πn) = 0 for n = 1, 2, 3, ... to evaluate
Eqs. (6.11) and (6.14). Substituting the values for an and bn back into Eq. (6.1), we obtain
the Fourier series for f (x):
∞ πnx
X 12
f (x) = (−1)n+1 sin (6.15)
n=1
nπ 2
12 πx 1 1 3πx
!
= sin − sin πx + sin − ... . (6.16)
π 2 2 3 2
Figure 6.3 shows f (x) together with a series of plots that include successively more terms in
the Fourier series for f (x). As more terms are included in the sum, the series approximates
the periodic function better and better. In fact, including only those terms up through n = 4
results in a surprisingly good approximation of the sawtooth waveform.
The sawtooth waveform shown in Fig. 6.3 is an odd function. That is f (x) = − f (−x).
Therefore, the Fourier series for f (x) must be odd, term by term. This means that no even
term, that is, no term proportional to cos kn x can contribute to the sum since cosine is an
171 Fourier series
t
(a) 6 (b)
x 0 x
−8 −6 −4 −2 2 4 6 8 1 2 3 4 5 6 7 8
−6 −6
t
Fig. 6.4 Fourier components that make up the Fourier series plotted in Fig. 6.3: (a) sine waves
with correct amplitudes (and phases) (b) amplitudes bn = (−1)n+1 12/nπ of sine waves.
even function and even one such term would spoil the odd symmetry of f (x). Therefore
we can set the coefficients an of the even cosine terms in Eq. (6.1) to zero. By noting
this symmetry, we can avoid having to explicitly perform the integrals in Eq. (6.11) to
determine the coefficients an . We only need to determine the coefficients bn of odd terms,
sin kn x, as only they can contribute to the Fourier sum.
Figure 6.4(a) shows the sine waveforms that were added together to make up the saw-
tooth wave. Note that they consist of the fundamental periodicity L and all its harmonics
L/n. Figure 6.4(b) shows the amplitudes bn of the sine waves that make up the sawtooth
wave. Note that because the power in a wave is proportional to the amplitude squared, the
power in each frequency is proportional to b2n . More generally, the power in each frequency
is proportional to a2n + b2n , but in this (special) case it reduces to b2n because an = 0 for all
n. We develop the subject of the power in a wave described by a Fourier series in §6.1.4.
The sets of coefficients {an } and {bn } give a quantitative measure of the harmonic content
of a particular waveform. They tell how much of a waveform is made up of sine and cosine
waves of a particular wavenumber kn , or equivalently of a wavelength λn = 2π/kn . As
noted in the introduction to this chapter, it is the harmonic content, that is the particular
collection of wavenumbers kn and their amplitudes an and bn , that gives a particular musical
instrument its characteristic sound.
To reinforce these ideas we look at another example. Consider the triangle waveform in
Fig. 6.1(b):
4x + 2 , −1 ≤ x ≤ 0 ,
(
f (x) = (6.17)
−4x + 2 , 0≤x≤1.
Note that it is an even function: f (−x) = f (x). Because it is an even function, the coeffi-
cients of the sin kn x terms are zero, i.e. bn = 0. Therefore, we only need to determine the
coefficients an , which are given by Eq. (6.6):
2 L/2
Z
an = f (x) cos kn x dx
L −L/2
Z 0 Z 1
2πnx 2πnx
! !
= (4x + 2) cos dx + (−4x + 2) cos dx
−1 2 0 2
16/(n2 π2 ) , if n is odd;
(
= (6.18)
0 , if n is even;
t
The resulting Fourier series for the triangle wave is
∞
X 16
f (x) = 2 π2
cos(nπx) (6.19)
n=1,3,5,...
n
16 cos(3πx) cos(5πx)
" #
= 2 cos(πx) + + + ... . (6.20)
π 9 25
6.1.2 Fourier series in time
The Fourier series developed above to describe periodic waveforms in space can also be
used to describe periodic waveforms in time. Doing so is as simple as replacing the spatial
variable x with time t and the wavenumber kn with the angular frequency ωn where
2πn
,
ωn = (6.21)
T
and T is the period of the waveform. With these changes, we write
∞ ∞
1 X X
f (t) = a0 + an cos ωn t + bn sin ωn t . (6.22)
2 n=1 n=1
In this case the Fourier series is a sum over sines and cosines of different temporal fre-
quencies ωn . Because of this correspondence between spatial and temporal Fourier series,
people often refer to the wavenumber kn that appears in the spatial Fourier series as a
spatial frequency, in analogy to the temporal frequency ωn . Thus, spatial frequency and
wavenumber are simply two names for kn .
6.1.3 Basis functions and orthogonality
The sine and cosine functions that we use in Fourier series are often referred to as ba-
sis functions. The integrals over their products, like those given in Eqs. (6.3) and (6.4),
are known as scalar products and express the orthogonality the the sine and cosine basis
functions.
This language of basis functions and orthogonality is an extension of the concepts and
language we use when speaking of a Cartesian coordinate system. In a Cartesian coordinate
system, for example, we can express the coordinate of any point a as linear combination
of the three orthogonal basis vectors, ê x , êy , and êz , the unit vectors along the x, y, and z
directions, respectively. Similarly, we can express any periodic function with period L as
linear combinations of the sin kn x and cos kn x basis functions (where kn = 2πn/L). Two
Cartesian vectors a and b are orthogonal if their dot product is zero, that is if
3
X
a·b= ai bi = 0 , (6.23)
i=1
where the subscripts i = 1, 2, 3 refer to the x, y, z components of each vector. We can

readily adapt the definition of the dot (or scalar) product to 100 component vectors (or any
other number). Consider, for example the two 100 component vectors A = cos 2πi/100 and
173 Fourier series
t
B = sin 2πi/100, where i runs from 1 to 100. Generalizing the dot product written above
to 100-dimensional vectors, we write the dot product of these two vectors as
100
X
A·B= Ai Bi . (6.24)
i=1
We can increase the number of components to be as large as we wish and the scalar product
written above still makes sense. In the case of Fourier series, we have simply increased
the number of of components to infinity, because there is an infinite number of values of
x between −L/2 and L/2. In so doing the sum over the various components in Eq. (6.24)
becomes an integral. For the Fourier series we have defined, there is also an infinite number
of basis functions cos kn x and sin kn x (where kn = 2πn/L) because n runs from 1 to infinity.
With these considerations in mind, the scalar products of the sine and cosine basis functions
are:
L/2
L
Z
cos kn x cos km x dx = δnm (6.25)
−L/2 2
Z L/2
L
sin kn x sin km x dx = δnm (6.26)
−L/2 2
Z L/2
sin kn x cos km x dx = 0 (6.27)
−L/2
If the scalar product, thus defined, is zero, then we say that the sine and cosine basis vec-
tors are orthogonal. Equations (6.25)-(6.27) show that all the different basis Fourier basis
functions are indeed orthogonal, as their scalar products are always zero. Only when per-
forming a scalar product of a basis vector with itself does the scalar product yield a finite
answer, similar to the behavior of the Cartesian basis vectors ê x , êy , and êz .
Rereading the first paragraph of this section, you should now understand the terminology
introduced there about basis functions, scalar products, and orthogonality.
Exercise 6.1.1 Write a computer program to compute the sum in Eq. (6.24) for A =
cos 2πi/100 and B = sin 2πi/100 where i runs from 1 to 100. Do you get zero?
Repeat the calculation for A·A and B·B. Do the numerical values of the dot products
you obtain make sense? Briefly explain.
6.1.4 Power in periodic waves
The power in a waveform is proportional to the square of the wave. For a periodic function
f (x) represented by a Fourier series, the energy per period (wavelength) is proportional to
t
the square of f (x) averaged over one period
2 L/2
Z
P∝ | f (x)|2 dx (6.28)
L −L/2
2
2 L/2 1
Z X∞ X ∞
= a0 + an cos kn x + bn sin kn x dx (6.29)
L −L/2 2 n=1 n=1

 Z L/2 ∞ L/2 ∞ Z L/2 
2  1 X Z X
2 2 2 2 2
=  a0 dx + cos kn x dx +

an bn sin kn x dx (6.30)
L 4 −L/2 n=1 −L/2 n=1 −L/2
∞ ∞
 
2  1 1 X 1 X
=  La20 + L a2n + L b2n 

(6.31)
L 4 2 n=1 2 n=1
a20 X∞
a2n + b2n .

= + (6.32)
2 n=1
The step from Eq. (6.29) to Eq. (6.30) follows from the orthogonality of the sine and
cosine functions, which makes all the integrals on the interval [−L/2, L/2] not involving
the square of a sine or cosine function equal to zero. The integrals over cos2 kn x and sin2 kn x
are simply L/2, the average value of cos2 kn x and sin2 kn x, which is 1/2, times the length
of the integration interval, which is L. The equality between Eqs. (6.28) and (6.32),
2 L/2 a2 X ∞
Z
| f (x)|2 dx = 0 + a2n + b2n

(6.33)
L −L/2 2 n=1
is known as Parseval’s theorem. In this context it means that the power Pn in the nth Fourier
component is proportional to the sum of the square of the Fourier amplitudes. For mathe-
matical convenience we make the proportionality in Eq. (6.28) an equality and define the
power spectrum as
Pn = |an |2 + |bn |2 , (6.34)
which gives the power in the nth Fourier component for n ≥ 1. Because the n = 0 term
is not associated with any oscillations, we generally do not include it in the power spec-
trum. Thus, the power contained in the part of the frequency spectrum characterized by the
wavevector kn is given by Pn , aside from a numerical constant that depends on the type of
wave being described—sound, electrical, etc. The total power per period is given by sum-
ming Pn over all n. For example, the power in the sawtooth wave depicted in Fig. 6.1(a)
with a Fourier series given by Eq. (6.19) is
∞ !2 ∞
X 12 144 X 1
Ptot = = 2 = 24 . (6.35)
n=1
nπ π n=1 n2
Performing the same calculation for the triangle wave depicted in Fig. 6.1(b) gives
∞ !2 ∞
X 16 256 X 1 8
Ptot = 2 π2
= = . (6.36)
n=1,3,5,...
n π4
n=1,3,5,...
n 4 3
175 Fourier series
t
Exercise 6.1.2 What are the wavelengths of the first three spatial frequencies kn with
non-zero power in the sawtooth and triangle waves depicted in Fig. 6.1? What frac-
tion of the total power in each waveform is contained in (i) the longest wavelength?
(ii) the longest two wavelengths with non-zero power? (iii) the longest three wave-
lengths with non-zero power?
6.1.5 Complex exponential Fourier series
Fourier series can also be expressed in exponential form by exploiting the identity eiθ =
cos θ + i sin θ. In this case, we write the Fourier series as
X∞
f (x) = An eikn x , (6.37)
n=−∞
where once again kn = 2πn/L and L is the period of the waveform. In general, the coeffi-
cients An are complex. Complex coefficients are required to correctly capture the phase and
parity (oddness or evenness) of f (x). The coefficients can be determined using a procedure
similar to the one we used to obtain the coefficients of the sine and cosine Fourier series.
In this case, we multiply both sides of Eq. (6.37) by e−ikm x and integrate over one period:
Z L/2 Z L/2  X ∞

−ikm x
dx = An e  e−ikm x dx
ikn x 

f (x) e


−L/2 −L/2 n=−∞
∞
X Z L/2
= An ei(kn −km )x dx
n=−∞ −L/2
∞
X eiπ(n−m) − e−iπ(n−m)
= An L
n=−∞
2iπ(n − m)
∞
X sin[π(n − m)]
= An L
n=−∞
π(n − m)
X∞
= An L δnm
n=−∞
= Am L , (6.38)
where we have used the fact that sin[π(n − m)] = 0 for n , m (n and m are integers) and
the fact that limξ→0 sin ξ/ξ = 1 (here ξ = π(n − m)).1 Thus, Eq. (6.38) gives
1 L/2
Z
An = f (x) e−ikn x dx (6.39)
L −L/2
It is left as an exercise for the reader to show that the coefficients An can be expressed in
terms of those for the sine and cosine series
 1
(a|n| + ib|n| ) , n < 0 ;
 12



An =  a0 , n=0; (6.40)

 21

(a ib ) , n > 0 .


2 n − n
1 You can show that limξ→0 sin ξ/ξ = 1 by writing the Taylor series
sin ξ/ξ = (ξ − ξ3 /3! + ξ5 /5! − ...)/ξ = 1 − ξ2 /3! + ξ4 /5! − ...), which for ξ → 0 is equal to 1.
t
Equations (6.37) and (6.39) constitute a completely equivalent alternative form to the sine
and cosine Fourier series expressed by Eq. (6.1) and Eqs. (6.6)-(6.8). The power spectrum
for waves described by a complex exponential Fourier series is given by
Pn = |An |2 = A∗n An . (6.41)
The primary importance of the exponential form of Fourier series is that the exponential
form is readily generalized to describe pulses—that is, functions of finite extent. This is
the subject of Fourier transforms, which are discussed in the next section.
6.2 Fourier transforms
When we encounter waves in nature, they have a finite spatial extent (and a finite temporal
duration). They are not infinitely periodic extending from x = −∞ to x = ∞, as is the case
for the waves described by the Fourier series encountered in §6.1. Nevertheless, it is still
useful to have some means of describing the harmonic content of such waves. As we shall
see, the mere fact that real waves have a finite extent or a finite duration means that more
than one frequency is required to describe their harmonic content.
6.2.1 Fourier transform pairs
We begin by considering a wave pulse that is simply a sinusoidal wave modulated by a

Gaussian envelope:
2
/2a2
f (x) = B e−x cos k0 x . (6.42)
Figure 6.5(a) shows a plot of Eq. (6.42) for the case that there are many oscillations in the
pulse, i.e. a k0−1 . The periodicity (wavelength) is approximately 2π/k0 = 2π/3 ' 2.1.
But unlike a cosine wave, this wave pulse has a finite duration of about 5a = 40. We cannot
use a Fourier series to evaluate the harmonic content of such a wave pulse because it does
not repeat. We could, however, use a Fourier series to evaluate the harmonic content of the
periodic wave in Fig. 6.5(b). It has been made periodic by repeating the wave pulse in Fig.
6.5(a) at intervals of x = 100. The exponential form of the Fourier series for this wave is
given by Eq. (6.37) with the Fourier coefficients given by Eq. (6.39) with L = 100 and f (x)
given by Eq. (6.42).
If our goal is to describe the harmonic content of a pulse such as Eq. (6.42), the value of
L we choose is arbitrary so long as we choose L to be much greater than the extent of the
wave pulse, i.e. so long as L a. As we make L larger and larger, our approximation of
f (x) should become better and better. In the general case, we let L go to infinity so that we
can describe a wave pulse of arbitrary (but finite) duration. To this end, we rewrite the sum
177 Fourier transforms
t
(a) f (x)
x
−20 −10 10 20
(b) f (x)
x
−200 −100 100 200
t
Fig. 6.5 (a) Cosine wave with Gaussian envelope: k0 = 3 and a = 8 (arrow in background
extends from −a to a). (b) The same cosine wave with Gaussian envelope as in (a)
but repeated at intervals of x = 100.
in Eq. (6.37) as
∞
X
f (x) = lim An eikn x (6.43)
L→∞
n=−∞
∞ L 2π
X !
= lim An ei 2πnx/L ∆n
L→∞
n=−∞
2π L
∞
X LAn i 2πnx/L 2π∆n
= lim e
L→∞
n=−∞
2π L
1 ∞
Z
= F(k) eikx dk , (6.44)
2π −∞
where limL→∞ 2π∆n/L → dk (note that ∆n = 1 so the sum is not changed by its intro-
duction) and we have defined the function F(k) = limL→∞ LAn . Just as the amplitude An
determines how strongly waves of spatial frequency kn = 2πn/L contribute the the Fourier
sum of Eq. (6.43), so the function F(k) determines how strongly waves of spatial frequency
k contribute the the Fourier integral of Eq. (6.44). Using the expression for An given by Eq.
(6.39), we obtain an expression for F(k):
1 L/2
Z !
F(k) = lim LAn = lim L f (x) e−ikn x dx (6.45)
L→∞ L→∞ L −L/2
Z ∞
= f (x) e−ikx dx . (6.46)
−∞
The function F(k) given by Eq. (6.46) is the Fourier transform of f (x). The inverse trans-
form of F(k) is given by Eq. (6.46). We rewrite this Fourier transform pair again for refer-
t
ence:
∞
dk ikx
Z
f (x) = e F(k) (6.47)
2π
Z−∞
∞
F(k) = dx e−ikx f (x) . (6.48)
−∞
For notational purposes, it is useful to define the Fourier transform operator and its inverse
Z ∞
F {...} ≡ dx e−ikx {...} (6.49)
−∞
Z ∞
F −1 {...} ≡ dx eikx {...} . (6.50)
−∞
Thus, we can write F { f (x)}, denoting the Fourier transform of f (x), which is equal to F(k).
Similiarly, we can write the inverse Fourier transform as F −1 {F(k)}, which is equal to f (x).
We shall make use of this notation in §6.2.4.
Equation (6.47) is the Fourier transform of f (x), which is a function of the spatial vari-
able x. Often we want to take the Fourier transform of a function g(t) that is a function
of time t. Mathematically, it is no problem to simply replace the spatial variable x with
the temporal variable t in Eqs. (6.47) and (6.48). The Fourier transform variable k then
becomes the angular frequency ω. Everything works fine if we make the simple substitu-
tion of t for x and of ω for k. However, in physics it is conventional to write the temporal
Fourier transform pairs with the opposite sign convention to that used in Eqs. (6.47) and
(6.48) above. That is, we write
dω
Z ∞
g(t) = G(ω) e−iωt (6.51)
−∞ 2π
Z ∞
G(ω) = dt g(t) eiωt . (6.52)
−∞
This choice of sign convention presents no difficulty mathematically, but it might seem odd
to use different sign conventions for spatial and temporal variables. The reason for doing
so to do with the fact that a wave propagating in the positive x direction is given by ei(kx−ωt) ,
that is, with a positive sign in front of the kx term and a minus sign in front of the ωt term.
This is the origin of the sign convention for spatial and temporal Fourier transforms often
used in physics, and we adopt it here.
We note that the Fourier transform pairs given by Eqs. (6.47) and (6.48) and by Eqs.
(6.51) and (6.52) do not incorporate the factor of 2π symmetrically. On the other hand, if
we rewrite the temporal Fourier transform pairs, in terms of the conventional frequency
ν = ω/2π, Eqs. (6.51) and (6.52) become
Z ∞
g(t) = dν e−i 2πνt G(ν) (6.53)
−∞
Z ∞
G(ν) = dt ei 2πνt g(t) . (6.54)
−∞
Written in terms of the frequency ν instead of the angular frequency ω, the temporal Fourier
transform pairs appear as completely symmetric transforms of each other. Remembering
this fact can help you remember which transform includes the factor of 2π.
t
F(k)
k
−4 −3 −2 −1 0 1 2 3 4
t
Fig. 6.6 Fourier transform of a Gaussian wave packet. The Gaussian wave packet is given by
Eq. (6.42) and is plotted in Fig. 6.5(a). Its Fourier transform, plotted above, is given by
Eq. (6.55).
6.2.2 Some examples
Having introduced the formalism of Fourier transforms, let’s use it to find the Fourier
transforms of a few functions. We start with Eq. (6.42), a Gaussian modulated cosine wave.
We use Eq. (6.48) to calculate the Fourier transform of Eq. (6.42):
Z ∞
F(k) = f (x) e−ikx dx
Z−∞
∞
2 2
= B e−x /2a cos k0 x e−ikx dx
−∞
π − 1 (k−k0 )2 a2
r
1 2 2
+ e− 2 (k+k0 ) a

= Ba e 2 (6.55)
2
The result, plotted in Fig. 6.6, is just what you might have expected: there are peaks in
the Fourier transform at k = ±k0 corresponding to the primary frequency of the cosine
wave. You may wonder why there are two peaks, one centered about k0 and the other
centered about −k0 . The reason is readily understood by recalling the complex exponential
representation of cosine:
cos k0 x = 12 eik0 x + e−ik0 x

(6.56)
A straightforward calculation of the Fourier transforms of eik0 x and e−ik0 x reveals the the
1 2 2 1 2 2
former is proportional to e− 2 (k−k0 ) a while the latter is proportional to e− 2 (k+k0 ) a . Since
cos k0 x contains both eik0 x and e−ik0 x , its Fourier transform has peaks centered at both k0
and −k0 .
The two peaks are Gaussian functions centered about k = ±k0 having a frequency width
of approximately
√ a−1 . To be more precise, the half width of the peaks at 1/e of the peak
height is 2/a. This means that the range of wavenumbers ∆k required to form a pulse of
length ∆x ' a is approximately the reciprocal of the pulse duration a−1 .
In the limit that k0 goes to zero, cos k0 t approaches unity so that the pulse shown Fig.
6.5 becomes a simple Gaussian without any oscillations. That is, in this limit
B 2 2
f (x) = √ e−x /2a , (6.57)
2π
t
and the Fourier transform becomes
1 2 2
F(k) = Ba e− 2 k a
. (6.58)
Thus, we see that the Fourier transform of a Gaussian function with a width of ∆x ' a is
another Gaussian function with a wavenumber width of ∆k ' a−1 . That is, the characteristic
widths of a function and its Fourier transform are inversely proportional:
1
∆x ∝ . (6.59)
∆k
Please note that the proportionality constant in the above relation is not the same for all
Fourier transform pairs but depends on the functional form of f (x). We will have more to
say about this inverse proportionality in the next section, where we make the statement in
a more general and mathematically precise way. For now, suffice it to say that it is a very
general and an extremely important property of pulses and their Fourier transforms. It is
known as the uncertainty relation or uncertainty principle for waves.
Another common function is the exponential:
g1 (x) = e−|x|/a . (6.60)
With a little algebra you can show that its Fourier transform is given by
2a
Z ∞
G1 (k) = e−|x|/a e−ikx dx = (6.61)
−∞ 1 + (ka)2
Similarly, the Fourier transform of
g2 (x) = e−|x|/a cos k0 x , (6.62)
is
a a
G2 (k) = + (6.63)
1 + (k + k0 )2 a2 1 + (k − k0 )2 a2
The functions G1 (k) and G2 (k) are Lorentzian functions centered about wavenumbers of
zero and ±k0 , respectively. Thus, the Fourier transform of an exponential is a Lorentzian.
Exercise 6.2.1 Perform the integral in Eq. (6.61) to calculate the Fourier transform of
g1 (x) = e−|x|/a . Hint: Start by explaining why the following equality is true:
Z ∞ "Z ∞ #
e−|x|/a e−ikx dx = 2 Re e−x/a e−ikx dx (6.64)
−∞ 0
Let’s do another example. Consider a square pulse defined by the equation
1 , |x| ≤ a
(
g3 (x) = (6.65)
0 , |x| > a ,
t
(a) g3 (x)
1
0.5
x
−3 −2 −1 0 1 2 3
(b) 0.8 G3 (k)
0.4
k
−20 −15 −10 −5 5 10 15 20
t
Fig. 6.7 (a) Square pulse and (b) its Fourier transform.
as shown in Fig. 6.7(a). The Fourier transform is given by

Z ∞
G3 (k) = g3 (x) e−ikx dx
Z−∞
a
= e−ikx dx
−a
a
e−ikx eika − e−ika 2 sin ka
= = =
−ik −a ik k
= 2a sinc ka , (6.66)
where we have used the definition sinc x ≡ sin x/x. Equation (6.66) is plotted in Fig. 6.7(b).
Note that since sin x and x are both odd functions, their quotient sin x/x is even.
Unlike the Fourier transforms of the Gaussian and Lorentzian, which decay monotoni-
cally with increasing k, the Fourier transform of the square pulse oscillates about zero as
it decays away. The origin of the oscillations—called “ringing"—is the abrupt change in
amplitude of the square pulse. Very high spatial frequencies k are required in the Fourier
transform F(k) in order to reproduce these sharp features. A useful rule of thumb is that
smooth features in a pulse produce smooth Fourier transforms while sharp features in a
pulse produce Fourier transforms with oscillations at large k.
6.2.3 The uncertainty principle
In the previous section we pointed out that the characteristic width of a pulse and its Fourier
transform are inversely proportional to each other. This idea is summarized by Eq. (6.59).
We would like to make the statement about the relationship between the width ∆x of a
pulse and its Fourier transform ∆k more precise. To do so we need a precise definition of
width of a function.
t
We start by defining a function P f (x)
| f (x)|2
P f (x) = R ∞ . (6.67)
−∞
| f (x)|2 dx
As the power in the waveform f (x) is proportional to | f (x)|2 , P f (x) dx gives the fractional
power in the wave between the positions x and x + dx. We then define the average position
of the wave as
Z ∞
hxi = x P f (x) dx , (6.68)
−∞
where the angular brackets h...i to denote the average. Thus, we define the average position
weighted according to what fraction of the wave’s power is located near any given point.
Similary, we can define the mean square width or variance of x as
Z ∞
Var(x) = (x − hxi)2 |P f (x)|2 dx . (6.69)
−∞
We define the width of a wave function as

∆x = Var(x) ,
p
(6.70)
or
Z ∞
2 2
(∆x) = h(x − hxi) i = (x − hxi)2 |P f (x)|2 dx . (6.71)
−∞
The quantity ∆x gives us a quantitative measure of how far in space a pulse extends from
its average position hxi. A similar equation applies to ∆k, as is illustrated in the worked
examples below.
Exercise 6.2.2 Using the definitions above, show that

h(x − hxi)2 i = hx2 i − hxi2 (6.72)
Let’s use Eqs. (6.71) and (6.72) to find ∆x and ∆k for some of the Fourier transform
pairs we considered above. We start with the Gaussian pulse defined by Eq. (6.57) and its
Fourier transform Eq. (6.58). We calculate (∆x)2 using Eq. (6.72), noting first of all that
hxi = 0 because f (x) is centered about x = 0. Therefore, in this case
(∆x)2 = hx2 i − hxi2 = hx2 i (6.73)
Z ∞
= x2 P f (x) dx (6.74)
−∞
2
1 B −x2 /2a2
Z ∞
2
= R∞ B x √ e dx (6.75)
−∞
| √2π e−x2 /2a2 |2 dx −∞ 2π
= 12 a2 . (6.76)
Thus we find
a
∆x = √ (6.77)
2
t
We can calculate ∆k following a similar procedure. First we see that hki = 0 because F(k)
is symmetric about x = 0. Therefore
(∆k)2 = hk2 i − hki2 = hk2 i (6.78)

Z ∞
= k2 PF (k) dx (6.79)
−∞
1
Z ∞
= R∞ k2 |F(k)|2 dk (6.80)
2
|F(k)| dk −∞
−∞
1
Z ∞
1 2 2
= R∞ k2 |B a e− 2 k a )|2 dk (6.81)
− 21 k2 a2 2
−∞
|B a e )| dk −∞
1
= 2 , (6.82)
2a
Thus we find
1
∆k = √ (6.83)
2a
Combining the expressions for ∆x and ∆k we obtain
1
∆x ∆k = (6.84)
2
Equation (6.84) is the precise statement of the uncertainty principle for a Gaussian wave
packet. Similar expressions can be obtained for other wave packets, for example, the expo-
nential pulse given by Eq. (6.60) and its Fourier transform Eq. (6.61). However, for wave
packets other than the Gaussian wave packet, the numerical value of the produce ∆x ∆k
is generally greater than 1/2. The Gaussian wave packet yields the minimum value of the
product ∆x ∆k. Therefore, the most general expression of the uncertainty principle is writ-
ten as an inequality to cover all cases:
1
∆x ∆k ≥ . (6.85)
2
Exercise 6.2.3 Find ∆x and ∆k for the exponential pulse given by Eq. √ (6.60) and its
Fourier transform Eq. (6.61). Show that the product of ∆x ∆k = 1/ 2. Is this consis-
tent with the uncertainty relation Eq. (6.85)?
6.2.4 Some theorems
The Fourier transforms of the functions given in Eqs. (6.42) and (6.57) illustrate a more
general property of Fourier transforms. If F(k) is the Fourier transform of f (x), i.e. if
F { f (x)} = F(k), then it is straightforward to show that
F { f (x)eik0 x } = F(k − k0 ) . (6.86)
Equation (6.111) is known as the frequency shifting theorem. Rewriting cos k0 x in complex
exponential form, the frequency shifting theorem can be used to obtain Eq. (6.42) from Eq.
t
(6.57) or to obtain Eq. (6.63) from Eq. (6.61). We can obtain a similar relation relating a
spatially shifted function to its Fourier transform
F { f (x − x0 )} = F(k)e−i kx0 . (6.87)
Below, we state without proof a few other useful theorems about Fourier transforms:
( n
d f (x)
)
F = (ik)n F(k) (6.88)
dxn
dn F(k)
F {xn f (x)} = in (6.89)
dkn
F { f (x) ∗ g(x)} = F(k) G(k) , (6.90)
where f (x) ∗ g(x) denotes the convolution of f (x) with g(x)
Z ∞
f (x) ∗ g(x) ≡ f (x) g(x − ξ) dξ . (6.91)
−∞
The power spectrum is defined as

P(k) = |F(k)|2 = F(k) F ∗ (k) . (6.92)
The power between spatial frequencies k and k + dk in the spectrum is equal to (aside from
some dimensional constant) P(k) dk. The total power is obtained by integrating P(k) dk
over all k.
6.3 Discrete Fourier Transforms
Fourier transforms provide a very powerful means for understanding and analyzing wave-
forms and other fluctuating functions. When using a computer or dealing with experimen-
tal data, the functions encountered are necessarily represented by discrete data sets defined
over a finite range rather than continuous functions defined over an infinite range. To ex-
tend the tools of Fourier transforms to such data sets, it is useful to develop a discrete
Fourier transform or DFT. In this case, the function that is to analyzed is represented by a
finite set of evenly-spaced data points a distance 4x apart, as illustrated in Fig. 6.8.
To see how the DFT works, we start by recalling the expressions for the continuous
Fourier transform pair:
Z ∞
F(k) = f (x) e−ikx dx (6.93)
−∞
1
Z ∞
f (x) = F(k) eikx dk . (6.94)
2π −∞
We wish to form the discrete equivalent of the continuous Fourier transform where the
function f (x) is uniformly sampled over a range of length L at discrete intervals ∆x = L/N,
where N is taken to be an even integer. The length L is chosen to correspond to the range
of x where f (x) is non-zero; ideally f (x) = 0 outside this range. To capture all the spectral
185 Discrete Fourier Transforms
t
information in f (x), the interval ∆x should be chosen so there are at least two points per
cycle of the highest spatial frequency present in the waveform.
In formulating the DFT, the convention is to choose the interval over which the function
f (x) is sampled to be 0 ≤ x ≤ L with
L
x = n∆x , n = 0, 1, 2, ..., N − 1 , ∆x = . (6.95)
N
Note that the interval is not symmetric about x = 0, in contrast to the definition of the
Fourier transform. This choice means has some practical implications for representing
waveforms with DFTs that we return to later. For the moment, we simply note that this
choice is the standard way that DFTs are defined.
Because there are only N independent samples of f (x), there are only N independent
values of F(k). The smallest absolute value of k is ∆k = 2π/L so that the discrete values of
k are given by
2π
k = l∆k , l = −N/2, ..., N/2 , ∆k = . (6.96)
L
The range of spatial frequencies required to discretely represent f (x) must still include
both negative and positive values. This follows from the definition of the Fourier transform,
which requires positive and negative spatial frequencies irrespective of the range of x over
which f (x) is non-zero.
f (x) (a)
f3
f2
f1 fN−2
f0 fN−1
x0 x1 x2 x3 x
xN−2 xN−1
(b)
n = N −1
(c)
n=1
t
Fig. 6.8 (a) Function f (x) sampled at evenly-spaced intervals for the discrete Fourier
transform (DFT). (b) Shortest wavelength (smallest spatial frequency) represented by
sampling at an interval ∆x. (c) Longest finite wavelength represented by DFT defined
over an interval of length N∆x.
t
Before forging forward with the mathematical development, we also note that the spatial
frequencies that are included in the DFT range from k = ∆k = 2π/L to k = ±(2π/L)(N/2) =
±2π/[L/(N/2)]. These two extreme values of wave vectors we are considering correspond
to (i) a wavelength equal to the entire interval L over which the data are defined (see Fig.
6.8(b)), and (ii) to a wavelength that is two data points long L/(N/2). This is so because
it takes a minimum of two data points to represent the oscillations of a sine wave, one
point for the negative trough and another point for the positive peak (see Fig. 6.8(c)). If
the function in question has contributions at higher frequencies, they will not be properly
captured by the DFT because the function is not sampled at sufficiently short intervals.
With these definitions, the discrete version of the product kx is given by
2π L 2π
! !
kx = (l∆k) (n∆x) = l n = ln (6.97)
L N N
Defining fn = f (n∆x), the discrete version of the Fourier transform becomes
Z ∞ N−1
X
F(k) = f (x) e−ikx dx ≈ F(l∆k) = fn e−i (2π/N)ln ∆x . (6.98)
−∞ n=0
The discrete Fourier transform (DFT) is defined as the last sum in Eq. (6.98) less ∆x:
N−1
X
Fl = fn e−i(2π/N)ln . (6.99)
n=0
With this definition of Fl , we can also write down the discrete version of the inverse Fourier
transform:
N/2−1
1 1
Z ∞ X
f (x) = F(k) eikx dk ≈ f (n∆x) = Fl ei (2π/N)ln , (6.100)
2π −∞ l=−N/2
N∆x
where dk/2π is rendered as

1 1 2π 1 2π 1
dk ≈ = = . (6.101)
2π 2π L 2π N∆x N∆x
The inverse discrete Fourier transform (iDFT) is defined as the last sum in Eq. (6.101) less
1/∆x:
N−1
1 X
fn = Fl ei(2π/N)ln . (6.102)
N l=0
Notice that in going from Eq. (6.101) to Eq. (6.102) we have changed the sum over l from
−N/2 to N/2−1 to 0 to N −1. There is no real problem with this change since Fl is periodic
in l with a period of N, as can be seen from Eq. (6.102). You just have to keep in mind that
the non-negative spatial frequencies appear in the first half (from l = 0 to l = N/2 − 1)
followed by the negative spatial frequencies (from l = N/2 to l = N − 1) in the second half,
as detailed in Table 6.1. The non-intuitive ordering of Fourier components in the DFT is
done, in part, so that the two sums in the definitions for the DFT and the iDFT, Eqs. (6.99)
and (6.102), respectively, both run from 0 to N − 1. Another reason for this choice is that
DFTs are useful primarily in computer applications where it is convenient for indices such
t
Table 6.1 Ordering of spatial frequencies for the discrete

Fourier transform.
positive frequencies negative frequencies

l kl l kl
0 0 N/2 (−N/2)∆k
1 ∆k N/2 + 1 (−N/2 + 1)∆k
2 2∆k N/2 + 2 (−N/2 + 2)∆k
.. .. .. ..
. . . .
N/2 − 2 (N/2 − 2)∆k N−2 −2∆k
N/2 − 1 (N/2 − 1)∆k N−1 −∆k
as n and l to run from 0 (or 1) to to some finite number, in this case N − 1 (or N). It is no
use protesting that you might prefer to do it another way. This is the standard that everyone
uses.
Discrete Fourier transforms can be calculated very efficiently using a computer algo-
rithm called the Fast Fourier Transform or the FFT. There are several variants of the FFT
algorithm, the oldest of which dates back to to Gauss in 1805, although the basic algo-
rithm has been independently rediscovered by a number of people since that time. Many
mathematical computer libraries implement the FFT algorithm. Those that do also gen-
erally provide a routine for calculating the discrete Fourier (spatial) frequencies kl that
match the corresponding Fl ; that is, the kl start from 0 and go to (N/2 − 1)∆k and then
from −(N/2)∆k to −∆k. Most FFT packages also have routines that reorder the frequencies
from lowest (most negative) to highest (most positive), and, respectively, the corresponding
Fourier components.2
Once you have obtained the DFT from an FFT algorithm or by direct computation, you
can make direct connection with the conventional continuous Fourier transform and its
inverse simply by supplying ∆x as follows:
N−1 
X
(2π/N)ln

F(l∆k) → Fl ∆x =   fn e−i  ∆x
 (6.103)
n=0
 N/2−1 
1  1 X
i (2π/N)ln
 1
f (n∆x) → fn =  Fl e  . (6.104)
∆x N l=−N/2 ∆x
Let’s work an example. Consider the the function
f (x) = e−|x−x0 |/a cos k0 x (6.105)
To be definite, we set a = 12 , k0 = 2π/λ0 with λ0 = 34 , and x0 = 4. We choose L = 2x0 = 8

1
and N = 512 so that ∆x = L/N = 64 . The wave packet given by Eq. (6.105) is plotted in Fig.
6.9 with the points at which fn = f (n∆x) indicated. We wish to find the DFT for the points
given by fn = f (n∆x). We use the FFT algorithm to calculate the DFT. The algorithm works
2 For example, in Python, the DFT of an array y is provided within the scipy.fftpack package as fft(y).
t
f (x) or f n
0.8 (a)
0.4
0.0 x
1 2 3 4 5 6 7 8
−0.4
−0.8
Fl
40 (b)
20
0 l
128 256 384 512
−20
−40
40 (c)
20
0 l
240 256 272 288
−20
−40
F(k)
(d)
0.4
k
−20 −10 10 20
−0.4
t
Fig. 6.9 Discrete Fourier transform of an exponential wave packet. (a) Exponential wave
packet given by Eq. (6.105). (b) Discrete Fourier transform vs. l. (c) Discrete Fourier
transform reordered so that negative spatial frequencies (l < 256) appear before the
positive spatial frequencies (l > 256). Dotted lines connect DFT points for clarity of
presentation. (d) The DFT multiplied by ∆x is plotted against the proper spatial
frequencies. The continuous solid line shows the analytical Fourier transform.
best if the number of points N in fn is a power of two. That is why we chose N = 512 = 29 .
The DFT is plotted in Fig. 6.9(b) with the positive spatial frequencies running from 0 to 255
and the negative spatial frequencies running from 256 to 511. The DFT is replotted in Fig.
6.9(c) but this time with the spatial frequencies increasing monotonically from negative
(n = 0 to 255) to positive (n = 256 to 511) spatial frequencies.3
The raw DFT is simply indexed by the (dimensionless) integers l, but often you want to
work with the DFT as a function of the proper spatial frequencies, with the proper units of
inverse length. This is a simple matter, you simply shift the indices l by N/2 so that l runs
from −N/2 to N/2 − 1 and then multiply by ∆k. Before performing this rescaling, you must
be sure that you have reordered the DFT as described in the previous paragraph so that the
3 As do many other computer languages, Python provides an FFT helper routine fftshift that shifts the order
of the FFT so that the negative frequencies come first, followed by the positive frequencies.
t
DFT is plotted starting with the negative spatial frequencies and ending with the positive
spatial frequencies.
Moreover, if you want the digital Fourier transform itself to have the same units as the
equivalent continuous Fourier transform would have, you must multiply the DFT by ∆x,
the spacing between the point of the original data, as indicated by Eq. (6.103). Figure
6.9(d) shows the DFT multiplied by ∆x and plotted against the proper dimensional spatial
frequencies ∆k. For comparison, Fig. 6.9(d) also shows the analytic Fourier transform of
f (x), which is given by
a a
!
F(k) = + eik0 x , (6.106)
1 + (k + k0 )2 a2 1 + (k − k0 )2 a2
where we have used Eq. (6.93) to calculate the Fourier transform. Note that because the
original function f (x) is centered about a finite value x0 = 4, there are oscillations, given
by the the exponential factor eik0 x . As expected, the DFT and the analytic Fourier transform
are consistent with each other.
Fourier analysis is extremely powerful and can be used in clever ways that may not
be immediately obvious. One of the more useful applications of Fourier transforms is for
filtering data. Let’s look at one illustrative example. Consider the noisy Gaussian waveform
f (x) shown in in Fig. 6.10(a). The DFT is shown in Fig. 6.10(b). What is not so apparent in
Fig. 6.10(b) is that there is some high spatial frequency components in the DFT due to the
noise in f (x). The DFT is shown with a magnified y-axis in Fig. 6.10(c), where the high
frequency noise is clearly visible. The noise in the high frequency components visible in
Fig. 6.10(c) and indicated by the gray arrows is removed in Fig. 6.10(d). Performing the
inverse DFT on the signal in in Fig. 6.10(d) returns the original Gaussian signal with the
noise removed.
You may have noticed that the DFT in Fig. 6.10(b) has oscillations in it. These arise
because the DFT assumes that f (x) goes from 0 to L, which means that all the x values are
positive. The DFT is effectively calculating the Fourier transform of f (x − x0 ) where x0 is
half the width of the data: x0 = 4. According to frequency shifting theorem Eq. (6.87), the
1 2 2
resultant Fourier transform is F(k)e−ikx0 , where in this case F(k) = e− 2 k a . Thus, the factor
of e−ikx0 introduces oscillations in the DFT. We could manually remove these oscillations
by multiplying the DFT by eikx0 . However, there is no need to do so as our purpose here is
only to filter out the high frequency components and then re-transform in order to recover
the original signal, less the noise. Once we obtain the Fourier filtered data, which is just a
list of fn values, we are free to plot those values against the original x data, which is what
we do in Fig. 6.10(d).
Finally we point out that the DFT is in general complex, with real and imaginary com-
ponents. In the examples above, the imaginary components are much smaller than the real
components. Nevertheless the imaginary components can be important, especially when
performing a DFT, manipulating it as we have done in the Fourier filtering example above,
and then performing the inverse DFT.
t
f (x)
(a)
1.0
0.5
x
−4 −3 −2 −1 0 1 2 3 4
F
100 l
(b)
50
0 l
64 128 192 256
−50
F
2 l
(c)
0 l
64 128 192 256
−2
F
2 l
(d)
0 l
64 128 192 256
−2
f (x)
(e)
1.0
0.5
x
−4 −3 −2 −1 0 1 2 3 4
t
Fig. 6.10 Filtering a noisy Gaussian pulse. (a) Gaussian pulse with noise. (b) DFT of Gaussian
pulse with spatial frequencies shifted so that they appear from negative to positive. (c)
DFT magnified along y direction so that noise is visible. Gray arrows indicate high
frequency noise. (d) Filtered DFT (magnified) with high frequencies removed (e)
Inverse DFT of filtered DFT plotted against x rather than DFT index.
Problems
6.1 Find the Fourier series for the following functions. Take note of the odd or even
symmetry of the functions to avoid calculating Fourier coefficients that must be zero
by symmetry. Sketch the function in each case.
191 Problems
t
(a)
x , if 0≤x≤1
(
f (x) = (6.107)
−x , if −1 ≤ x ≤ 0 .
(b)
0 , if 0≤x≤1
(
f (x) = (6.108)
−x , if −1 ≤ x ≤ 0 .
(c) Taking k0 to be a constant, what is the lowest spatial frequency in the Fourier
series for this function and how does it compare to k0 ?
− sin k0 x , if −π ≤ k0 x ≤ 0
(
f (x) = (6.109)
sin k0 x , if 0 ≤ k0 x ≤ π .
6.2 The following periodic function is defined on the interval [−2, 2] but repeats indefi-
nitely along the ±x axis:
−6(2 + x) , if −2 ≤ x ≤ − 32




f (x) =  2x , if |x| ≤ 23 (6.110)


3

 6(2 − x) , if


≤ x ≤ 2 .
2
(a) Explain how the Python function below codes for f (x). In particular, explain pre-
cisely how the mod works. Then use the function to plot the function over the in-
terval −5 < x < 5. Hint: look at the output of the expression x = np.mod(x+2.0,
4.0) - 2.0 for different values of x over the interval −5 < x < 5.
import numpy as np
def sawtooth(x):
x = np.mod(x+2.0, 4.0) - 2.0
y = np.where(np.abs(x) <= 1.5, 2.*x,
6.*(2.-np.abs(x))*np.sign(x))
return y
(b) Find the Fourier coefficients for the sine and cosine Fourier series. You will need
to perform the Fourier integral in a piecewise fashion splitting it up into different
sections over the integration interval.
On separate plots, use Python to plot the Fourier series keeping (i) the first non-
zero term in the Fourier series, (ii) the first 2 non-zero terms in the Fourier series,
(iii) the first 3 non-zero terms in the Fourier series, and (iv) the first 10 non-zero
terms in the Fourier series.
(c) Find the power spectrum of f (x), that is the power P(kn ) at each wavevector kn
and make a plot of P(kn ) vs. kn . Find the total power in the waveform by summing
the contributions from all Fourier components. Which Fourier component carries
the most energy?
(d) Find the Fourier coefficients for the exponential Fourier series using Eq. (6.39).
By including both positive and negative coefficients, show that the complex
Fourier series gives the same results as the sine and cosine series for the first
three non-zero terms in the Fourier Series.
t
6.3 Find the Fourier transforms of the following functions. Using Python, plot the func-
tions below and on a separate graph, plot both the real and imaginary parts (if they
exist) of the Fourier transforms of each function. From your plots, estimate ∆k (or
∆ω) and ∆x (or ∆t) and show that ∆k ∆x ∼ 1 (or ∆ω ∆t ∼ 1).
Take care about the absolute values in the definitions of f1 (x) and f2 (x) below
when performing the Fourier integrals.
(a) f1 (x) = 5e−3|x|
(b) f2 (x) =(5e−3|x| sin 2x
2(1 − 41 t2 ) , if |t| ≤ 2
(c) g(t) =
0 , if |t| > 2 .
6.4 The (spatial) frequency shifting theorem says the if f (x) and F(k) are Fourier trans-
form pairs, i.e. f (x) ⇐⇒ F(k) [sometimes written as F { f (x)} = F(k)], then
f (x) eik0 x ⇐⇒ F(k − k0 ) (6.111)
(a) Prove the frequency shifting theorem. That is, given that F(k) is the Fourier
transform of f (x), show that Eq. (6.111) is true.
(b) In Exercises 6.3(a) and 6.3(b), you found the Fourier transforms of the functions
f1 (x) = 5e−3|x| and f2 (x) = 5e−3|x| sin 2x by explicitly calculating the Fourier inte-
grals for each function. For this exercise, you are to obtain the Fourier transform
of f2 (x) = 5e−3|x| sin 2x by applying the frequency shifting theorem to the Fourier
transform of f1 (x) = 5e−3|x| . This illustrates how the frequency shifting theorem
can make the calculation of the Fourier transforms of functions involving sines
and cosines significantly easier.
6.5 Use the FFT feature of SciPy to find the discrete Fourier transforms of the functions
listed in problem 6.3. Take care to adequately sample, but not drastically oversample
the three functions. The number of points you use should be a power of 2. Plot the
discrete Fourier transforms you calculate and compare them (properly shifted and
scaled) to the continuous Fourier transforms you determined analytically in problem
6.3. After finding the discrete Fourier transforms, perform the inverse transforms to
and show that you get back the original function in each case.
6.6 This problem explores the use of spatial filtering with FFTs to clean up noisy data.
Start by considering the following Python routine. It generates a noisy Guassian
function and then performs a discrete Fourier transform (DFT). It then zeros the
high frequency components (recall the nonintuitive ordering of the frequenies for
FFT routines), and then inverse transforms the filtered transform to return a clean
Gaussian signal.
from numpy.fft import fft, fftshift, ifft
import numpy as np
import matplotlib.pyplot as plt
def gaussNoisy(x, noiseAmp):

return np.exp(-x*x/2.0) * (1.0 + noiseAmp*(np.random.randn(len(x))))
193 Problems
t
N = 256
x = np.linspace(-4.0, 4.0, N)
y = gaussNoisy(x, 0.1)
yFT = fft(y) # DFT in usual order with + frequencies first

YFT = fftshift(yFT) # DFT reordered so that - frequencies are first
yFT[6:N-6] = 0.0+0.0j # Zero high frequency components of DFT
yFTFT = ifft(yFT) # iDFT of cleaned up DFT
plt.figure(1, figsize=(8, 10) )
plt.subplot(511)
plt.plot(x,y) # Noisy Gaussian function
plt.subplot(512)
plt.plot(YFT.imag) # Imaginary part of DFT
plt.subplot(513)
plt.plot(YFT.real) # Real part of DFT
plt.subplot(514)
plt.plot(fftshift(yFT)) # Cleaned up DFT
plt.subplot(515)
plt.plot(yFTFT) # Inverse of cleaned up DFT
plt.show()
(a) Set the noise amplitude to zero in the and find the DFT for the clean Gaussian
function. How far away from the maximum in the DFT do you have to go for all
the DFT amplitudes to fall below 1% of the maximum.
(b) Generate the proper Fourier transform variable k for the Fourier transformed data
and replot F(k) vs. k following the example of Fig. 6.9(d).
(c) Rescale the y-axis for the plot of the DFT so that the high frequency Fourier
components due to the noise are clearly visible, much as we did in Fig. 6.9(c).
(d) In the program above, all the DFT components more than 6 points from the
center spatial frequency are zeroed. What value of k does this correspond to? Is
this a good choice or is there a better choice? Justify your answer.
(e) Why is the imaginary component of the FFT so much smaller than the real com-
ponent? Is there some signal there or is it all noise?
(f) Vary the amplitude of the noise set by the noiseAmp variable. Currently the
t
noise amplitude is set to 0.3 or 30% of the signal. How well does the Fourier
filtering work with the noise is significantly larger, say 100%?
195 Problems
t
x
7 Strings
In the previous chapter we saw how propagating and standing waves can arise in systems
of discrete masses coupled by springs or tethered to each other by an idealized massless
string. In this chapter we begin our exploration of waves in continuous media, that is,
in materials like the strings of a violin or guitar. Of course, strings are made of discrete
masses—molecules—but when the wavelength is much greater than the typical distance
between masses, it is generally more efficient to regard the mass a being continuously
distributed along the string rather than keeping track of the individual particles. Mathe-
matically, this means that instead of describing the system in terms of a set of N coupled
ordinary differential equations with one independent variable, time, we describe the system
in terms of a single partial differential equation with two independent variables, time and
distance along the string. The study of waves on strings in terms of a partial differential
equation in time and space sets the stage for describing a host of other kinds of waves,
including sound waves in gasses and liquids, surface waves on membranes or liquids (e.g.
drum heads and soap films, or ripples on a lake), and even electromagnetic waves (e.g.
radio or light).
7.1 Waves on strings
We begin by considering the same system we analyzed at the beginning of Chapter 5: a

system of N particles of mass m tethered to each other by a massless string stretched to
a tension T between two walls, as shown in Fig. 5.1. This will serve as our model for a
string with with mass per unit length of µ = m/l. The idea is to pass to the continuum limit
by simultaneously letting the distance l between masses to go to zero while each mass m
becomes proportionally smaller such that the ratio µ = m/l remains constant.
The equation of motion for an interior mass in the system is given by Eq. (5.2), which
we reproduce here for easy reference:
T
mÿ s = (y s−1 − 2y s + y s+1 ) , (7.1)
l
where s = 1, 2, . . . , N indexes the masses starting from left to right and y s is the trans-
verse displacement of mass s. The mass per unit length of the string is given by µ = m/l.
196
197 Waves on strings
t
Substituting m = µl into Eq. (7.1) and regrouping terms a bit gives
T y s−1 − 2y s + y s+1
!
ÿ s = , (7.2)
µ ∆x2
where we have rewritten l as ∆x to remind us that we intend to take the limit ∆x → 0 while
keeping µ constant (see Fig. 7.1). In the continuum description, the transverse displacement
y s (t) at a particular point along the string is specified by the distance x = s∆x: thus, y s (t) →
y(x, t). Similarly, the second time derivative at a point s, which is given by ÿ s (t), becomes
a partial derivative with respect to time with the spacial variable x held constant:
∂2 y(x, t)
ÿ s (t) → . (7.3)
∂t2
The terms in parentheses on the right of Eq. (7.2) are evaluated at a fixed time t where
y s (t) → y(x, t) , y s−1 (t) → y(x − ∆x, t) , y s+1 (t) → y(x + ∆x, t) . (7.4)
In the limit ∆x → 0, the term in parentheses in Eq. (7.2) becomes
y(x+∆x,t)−y(x,t)
y(x − ∆x, t) − 2y(x, t) + y(x + ∆x, t) ∆x − y(x,t)−y(x−∆x,t)
∆x
lim = lim
∆x→0 ∆x2 ∆x→0
∆x
∂y(x,t)
∂x x+∆x/2 − ∂y(x,t)
∂x x−∆x/2
= lim
∆x→0 ∆x
2
∂y (x, t)
= . (7.5)
∂x2
Substituting Eqs. (7.3) and (7.5) into Eq. (7.2) gives
∂2 y(x, t) T ∂2 y(x, t)
= . (7.6)
∂t2 µ ∂x2
The partial derivative on the left hand side of Eq. (7.6) gives the transverse acceleration of
t
Fig. 7.1 Top to bottom: The number of masses N increases as m decreases such that µ = m/l
remains constant to obtain a continuous mass distribution along string (∆x = l).
198 Strings
t
the string and has dimensions of [length]/[time2 ], while the partial derivative on the right
hand side has units of [length−1 ]. Thus, the term T /µ must have units of [length]2 /[time2 ]
or of [velocity]2 . In view of this we define a velocity
3 = T /µ ,
p
(7.7)
and rewrite Eq. (7.6) as
∂2 y(x, t) ∂2 y(x, t)
= 32 . (7.8)
∂t 2 ∂x2
Equation (7.8) is known as the wave equation. While developed here for a string with a
linear mass density µ held under tension T , the wave equation given by Eq. (7.8) turns
out to have much broader applicability. For example, it also applies to sound waves and
electromagnetic waves (e.g. light), where 3 turns is the velocity of the wave, whether it is
a string wave, a sound wave, or an electromagnetic wave.1 On the other hand, the wave
equation is not the only equation whose solutions can be waves and it does not apply to
waves in all circumstances. For example, surface waves in bodies of water are governed by
a different set of equations.
The wave equation is a restatement of Newton’s second law applied to a string. To em-
phasize this, we rewrite Eq. (7.6) as
∂2 y(x, t) ∂2 y(x, t)
µ = T . (7.9)
∂t2 ∂x2
The left side of Eq. (7.9) is simply mass (per unit length) times acceleration, while the
right side is the restoring force. The second spatial derivative of the transverse displace-
ment ∂2 y/∂x2 is the curvature of the string. The curvature gives the direction and the net
scaled restoring force: negative curvature at a point along the string means the force at that
point is in the −y direction, while positive curvature means the force is in the +y direction.
Examination of the sinusoidal wave sketched in Fig. 7.1 illustrates how this is consistent
with our expectations. For a sinusoidal disturbance, points above the equilibrium posi-
tion have a negative curvature, which means that those portions of the string experience
a force downwards towards the equilibrium zero-vertical-displacement position (y = 0).
Conversely, points below the equilibrium position have a positive curvature and experi-
ence a force upwards towards the equilibrium zero-vertical-displacement position.
7.2 Normal modes of a string
Having obtained the equation of motion, namely the wave equation, our first task is to
determine its normal modes. For the case of N discrete masses on a massless string, we
1 The wave equation for strings, which is based on Eqs. (5.2) and (7.1), is valid only in the limit that the
restoring force is linear in the displacement (see §1.1 and §5.1). Similarly, the more general wave equation is
valid for sound and electromagnetic waves only under certain circumstances, which generally correspond to
similar assumptions about the linearity of the restoring force.
199 Normal modes of a string
t
began with Eq. (5.11) by assuming a solution of the form
y s (t) = a s e−iωt , (7.10)
where the index s ranged from 1 to N. For a string, the mass is distributed continuously
along the string rather than in N discrete masses, and we let y s (t) → y(x, t). Following the
procedure developed in the previous section, we write the trial solution as
y(x, t) = a(x) e−iωt . (7.11)
Substituting Eq. (7.11) into Eq. (7.8), we obtain
d2 a(x) −iωt
−ω2 a(x) e−iωt = 32 e . (7.12)
dx2
The derivative of a(x) is a full derivative because a(x) is a function only of x. Canceling
the common oscillating exponentials and rearranging terms, this becomes
d2 a(x)
+ k2 a(x) = 0 , (7.13)
dx2
where we have defined
k = 3/ω . (7.14)
Equation (7.13) has the same familiar form as Eqs. (1.5) and (1.30) and thus, aside from a
change of variables, has the same oscillatory solutions, which can be expressed in terms of
complex exponentials
a(x) = Ae±ikx . (7.15)
Using Eq. (7.11), our solution becomes y(x, t) = Aei(±kx−ωt) . Noting that our trial solution
could have equally well started with eiωt , our general solution becomes
y(x, t) = Aei(±kx±ωt) (7.16)
which means there are four independent solutions
y(x, t) = A ei(kx−ωt) + A∗ e−i(kx−ωt) + B ei(kx+ωt) + B∗ e−i(kx+ωt) , (7.17)
where the constants A and B are complex and determined by the initial and boundary
conditions. The first and second and the third and fourth terms must be complex conjugates
of each other so that the solutions y(x, t) is real, as discussed in §1.5.2. As usual, the general
solution can be written in terms of sine and cosine functions rather than as the complex
exponentials used in Eq. (7.17). Which form is used is largely a matter of mathematical
convenience. In each case, however, there are four integration constants, consistent with
the wave equation being a second order partial differential equation with two independent
variables x and t. The four integration constants can be determined from the boundary and
initial conditions.
200 Strings
t
7.2.1 String with fixed ends
We wish to find the normal modes of a string held at a tension T by stretching its ends
between two fixed points, one at x = 0 and the other at x = L. In this case, it is convenient
to write the general solution in terms of sines and cosines
y(x, t) = A1 cos (kx − ωt) + A2 sin (kx − ωt) + B1 cos (kx + ωt) + B2 sin (kx + ωt) , (7.18)
where the amplitudes are real. The boundary conditions stipulate that y(0, t) = 0 and
y(L, t) = 0. Applying the boundary condition y(0, t) = 0 gives
y(x, t) = A1 cos ωt − A2 sin ωt + B1 cos ωt + B2 sin ωt

= (A1 + B1 ) cos ωt − (A2 − B2 ) sin ωt = 0 , (7.19)
which is satisfied for all t only if B1 = −A1 and B2 = A2 . Thus, Eq. (7.18) becomes
y(x, t) = A1 cos (kx − ωt) + A2 sin (kx − ωt) − A1 cos (kx + ωt) + A2 sin (kx + ωt)
= A1 [cos (kx − ωt) − cos (kx + ωt)] + A2 [sin (kx − ωt) + sin (kx + ωt)]
= 2A1 sin kx sin ωt + 2A2 sin kx cos ωt , (7.20)
where the final step follows from the usual sine and cosine addition formulas. Equation
(7.20) can be further simplified
y(x, t) = sin kx [2A1 sin ωt + 2A2 cos ωt]

= C sin kx cos (ωt − ϕ) . (7.21)
The boundary condition y(L, t) = 0 is satisfied at all times only if sin kL = 0, which means
that kL = nπ, where n = 1, 2, . . . , or
nπ
k= , n = 1, 2, . . . (7.22)
L
Substituting this result into Eq. (7.21), our solution for the normal modes is given by
nπx
y(x, t) = C sin cos (ωt − ϕ) , n = 1, 2, . . . (7.23)
L
The amplitude C and phase ϕ are usually determined by the initial conditions. Figure 7.2
shows the first four normal modes for t = 0 and ϕ = 0◦ and 180◦ . The normal modes
consist of sine waves with wavelengths λ given by
2π 2L
λ= = , (7.24)
k n
meaning that a half-integral number of wavelengths fit between the fixed endpoints a dis-
tance L apart.
The oscillation frequency is obtained from Eq. (7.14)
2L3 3
ω = k3 = = 2π , (7.25)
n λ
where for strings 3 is a constant fixed by Eq. (7.7) as 3 = T /µ.
p
The normal modes for a string look very similar to the normal modes for discrete masses
t
4
1 3
2
n=1
0 x
0.2 0.4 0.6 0.8 1.0
n=1
2
−1 3
4
t
Fig. 7.2 First four normal modes of a string. Dashed curves are 180◦ out of phase with solid
curves.
on a string, at least for the low-n long-wavelength modes. Compare the vibrations for the
n = 1 through 4 modes for a string shown in Fig. 7.2 to those for 16 alternating masses on
a string shown in Fig. 5.12. The modes look nearly the same. The reason is that when the
wavelength is much larger than the distance over which the mass along the string varies,
the mass distribution can be approximated as being very nearly continuous. In this case,
we say there is a separation of length scales: the length scale set by the wavelength and the
length scale set by the spatial variation in mass (i.e. the distance between masses) are very
different.
By contrast, when the wavelength becomes comparable to the length scale over which
the mass fluctuates, approximating the string as a continuous mass distribution no longer
makes sense. There are no modes for the string, for example, that look like the n = 15
modes for either the acoustic or optical branches of the alternating discrete masses shown
in Fig. 5.12, where the length scale of the mass fluctuations is essentially the same as the
wavelength.
The full solution to the problem of a string vibrating between two fixed ends can be
expressed as a sum over normal modes
∞
X
y(x, t) = Cn sin kn x cos (ωn t − ϕn ) , (7.26)
n=1
where from Eq. (7.22) kn = πn/L and from Eq. (7.25) ωn = kn 3 = 2L3/n where n = 1, 2, . . ..
The coefficients Cn and phases ϕn can be determined using methods similar to those used
to find the amplitudes of the terms in Fourier series discussed in §6.1.1. First we multpiply
both sides of Eq. (7.26) by sin km x with m = 1, 2, . . ., where km = πm/L, and integrate over
202 Strings
t
all x:
Z L ∞
X Z L
y(x, t) sin km x dx = Cn cos (ωn t − ϕn ) sin kn x sin km x dx , (7.27)
0 n=1 0
The integral on the right hand side obeys the usual orthogonality conditions
Z L
2
sin kn x sin km x dx = δn,m , (7.28)
0 L
where δn,m is the Kronecker-δ, which is equal to 1 if n = m and zero otherwise. The only
term in Eq. (7.27) that is nonzero is the n = m term. Evaluating Eq. (7.27) at t = 0 gives
2 L
Z
Cn cos ϕn = y(x, 0) sin kn x dx , (7.29)
L 0
where we have assumed we know the initial displacement y(x, 0) of the string. We have a
single equation for two unknowns Cn and ϕn , so we need another condition. We can obtain
another equation if we also assume we know the initial velocity ẏ(x, 0) along the length of
the string. Taking the partial derivative of Eq. (7.26) gives
∞
∂y(x, t) X
=− Cn sin kn x ωn cos (ωn t − ϕn ) . (7.30)

∂t n=1
As before, we multiply both sides by sin km x, integrate over x from 0 to L, and apply the
orthogonality condition. Evaluating the result at t = 0 gives
1 2 L
Z
Cn sin ϕn = ẏ(x, 0) sin kn x dx , (7.31)
ωn L 0
Equations (7.29) and (7.31) can then be combined to obtain Cn and ϕn .
As an example, consider a string of length L whose center point is displaced a distance
A transverse to the equilibrium position of the string. The equation for the initial position
of the string is

2Ax/L
 if 0 ≤ x ≤ L/2
y(x, 0) =  . (7.32)

2A(1 − x/L) if L/2 < x ≤ L/2

Equation (7.31) implies that ϕn = 0 for all n since the initial velocity ẏ(x, 0) is zero, which
leaves only Eq. (7.29) to evaluate. Because of the symmetry of the initial condition, only
the odd n terms in the sum in Eq. (7.29) are nonzero. Moreover, for the odd n terms, the
integral from 0 to L/2 is the same as the integral from L/2 to L.2 Thus, the expression for
the coefficients becomes
2 L/2 2Ax (n+1)/2

8A/(π2 n2 ) if n = 1, 3, . . .

πnx
Z
−(−1)

Cn = 2 sin dx =  . (7.33)

L 0 L L 0 if n = 2, 4, . . .

2 To see this, make a crude (pencil & paper) plot of y(x, 0) and of the first few sin kn x for odd n where
kn = πn/L.
t
Substituting Eq. (7.33) into Eq. (7.26) with ϕn = 0, the full solution is
∞
X (−1)(n+1)/2 nπx nπ3t
y(x, t) = −8A sin cos , (7.34)
n=1,3,5,...
πn2 2 L L
where we recalled that ωn = 3kn .
7.2.2 String with periodic boundary conditions
Next we examine the modes of a string for the case of periodic boundary conditions, just
as we did in the case of discrete masses. Here we imagine a string of length L held under
tension T with the periodic boundary conditions
y(x, t) = y(x + L, t) . (7.35)
You can think of periodic boundary conditions as meaning that the end of the string is tied
to the start of the string to form a continuous circle of circumference L. It’s not clear how
you would maintain tension in such a string, but it’s a convenient mathematical artifice as
it provides a convenient means of avoiding physical barriers.
Starting from Eq. (7.16), we apply this periodic boundary condition
Aei(±kx±ωt) = Aei(±k(x+L)±ωt) . (7.36)
Upon canceling common factors we obtain
1 = e±ikL , (7.37)
which is satisfied only when
2nπ
k= . (7.38)
L
Since k = 2π/λ, this condition corresponds to λ = L/n with n = 1, 2, . . ., which simply
means that to satisfy the periodic boundary conditions, an integral number of wavelengths
must fit into one circumference around the circle.
As for the case of a string with fixed ends, we can write the general solution for this
problem in terms of sines and cosines given by Eq. (7.18). The periodic boundary condi-
tions only restrict the allowed values of k, with the allowed values given by Eq. (7.38). The
values of A1 , A2 , B1 , and B2 are determined by the initial conditions.
7.2.3 Traveling waves
The general solution to the wave equation can be written in terms of sines and cosines,
as expressed in Eq. (7.18). Let’s consider just the first two terms in Eq. (7.18) and set
B1 = B2 = 0. The the solution takes the form
y(x, t) = A1 cos (kx − ωt) + A2 sin (kx − ωt) . (7.39)
This solution can also be written as
y(x, t) = A sin (kx − ωt − ϕ) . (7.40)
204 Strings
t
y x0 x1 x2
0.02
0.01
0.00 x
0.5 1.0 1. 5 2.0
− 0.01
t
− 0.02
Fig. 7.3 Traveling wave shown at three different different, equally-spaced times: t0 = 0 (solid),
t1 (long dashes), t2 = 2t1 (short dashes). Arrows indicate the positions x0 , x1 , and x2
of the first crest at times t0 , t1 , and t2 .
Let’s examine this equation. At a given instant in time, the curve maps out a sinusoidal
function in space. Figure 7.3 shows the spatial variation of y(x, t) at three different equally-
spaced times. From Fig. 7.3, we see that the sine wave appears to move to the right. How-
ever, it is important to be clear about what is moving, and where it is moving. If we focus
on a given piece of the string, say at x = x0 , then the motion of that piece of string is given
by
y(x0 , t) = A sin (kx0 − ωt − ϕ) . (7.41)
Equation (7.41) tells us that as time progresses, the piece of string at x = x0 simply moves
up and down in the transverse y direction; it does not move to the left or right. The gray
segments in Fig. 7.3 at x = x0 = 0.25 show that that piece of string moves downward from
being at the crest of the wave at when t = t0 = 0 to successively lower y positions at times
t = t1 and t = t2 = 2t1 .
On the other hand, different pieces of the string reach the maximum vertical displace-
ment at different times, and where that maximum occurs moves to the right as time pro-
gresses. Thus, we say that the disturbance—the wave—moves to the right even though in-
dividual pieces of the string move only up and down. Such a wave is known as a traveling
wave. We can figure out how fast the disturbance is moving by noting that the maximum
occurs when the argument of the sine function kx − ωt − ϕ, which we call the phase Φ, in
Eq. (7.40) is an odd multiple of π/2. In Fig. 7.3, we focus on the case of Φ = π/2, but the
analysis is the same for any other fixed value of Φ. Now we just need to determine how
fast that point of constant phase moves. The phase Φ = kx − ωt − ϕ is constant when its
derivative is zero:
dΦ d dx
= (kx − ωt − ϕ) = k −ω=0
dt dt dt
dx ω
⇒ = =3,
dt k
t
where the very last equality follows from Eq. (7.14). Thus we see that the derivative of the
point of constant phase moves to the right with a velocity of ω/k, which we defined earlier
as 3. This velocity
3 = ω/k , (7.42)
is known as the phase velocity. It is the velocity at which the the crests and nodes move.
We can write the argument of Eq. (7.39) as
ω
kx − ωt = k x − t (7.43)
k
= k(x − 3t) (7.44)
This means that we can write Eq. (7.39) as
y(x, t) = A1 cos [k(x − 3t)] + A2 sin [k(x − 3t)] . (7.45)
Our independent variables, x and t, always appear together in the combination x − 3t so that
we can write
y(x, t) = f (x − 3t) , (7.46)
where in this case f (ξ) = A1 cos kξ + A2 sin kξ. Alternatively, using the form given by Eq.
(7.40), we could write f (ξ) = A sin (kξ − ϕ). Either way, y(x, t) is a function of ξ = x − 3t.
When the independent variables x and t always appears in the combination ξ = x − 3t, this
tells us that the function f (x − 3t) moves to the right at a velocity 3.
We can apply the same analysis to the solution given by Eq. (7.18) for the case when the
coefficients A1 = A2 = 0, but B1 and B2 are not set to zero. However, in this case we see
that this part of the solution is a function of x + 3t. That is
y(x, t) = B1 cos (kx + ωt) + B2 sin (kx + ωt) (7.47)
= B1 cos [k(x + 3t)] + B2 sin [k(x + 3t)] (7.48)
= g(x + 3t) (7.49)
This solution corresponds to a sinusoidal wave moving to the left, that is, in the −x di-
rection. Thus, the full solution can be viewed as the linear superposition of right and left
propagating sinusoidal waves
y(x, t) = f (x − 3t) + g(x + 3t) . (7.50)
7.2.4 Standing waves
It is interesting to combine the solutions of left and right traveling waves. Consider the
combination when we set the amplitudes of the left and right traveling waves equal
y(x, t) = A[sin(kx − ωt) + sin(kx + ωt)] (7.51)
Applying the sine addition and subtraction formulas, this becomes
y(x, t) = 2A cos ωt sin kx . (7.52)
206 Strings
t
This is simply a sine wave in space, sin kx, with a time-varying amplitude of 2A cos ωt.
The crests and nodes are do not move to the left or right but remain stationary. Such a wave
is known as a standing wave. From Eq. (7.51) we see that consists of two identical but
counter-propagating sinusoidal waves. This is a general property of standing waves.
The normal modes of a string vibrating between two fixed ends are standing waves,
as you can readily verify by examining Eq. (7.23) and Fig. 7.2. As we shall see in the
next section, the counter-propagating sinusoidal waves that make up these standing wave
normal modes are the result of the reflection of waves at the boundaries, which create
identical waves traveling in opposite directions. The reflecting counter-propagating waves
add to form a standing wave normal mode.
7.3 General solutions of the wave equation
In the previous section we learned that the sinusoidal solutions to the wave equation, Eq.
(7.8), could be written in the form
y(x, t) = f (x − 3t) + g(x + 3t) , (7.53)
where f (ξ) and g(ξ) are sine and cosine functions. In fact, the functions f (ξ) and g(ξ) do
not need to be sines and cosines in order to satisfy the wave equation; they can be any
physically reasonable functions. We verify this by substituting y(x, t) = f (x − 3t) + g(x + 3t)
x
0.5 1.0 1.5 2.0
y
x
0.5 1.0 1.5 2.0
y
x
0.5 1.0 1.5 2.0
t
Fig. 7.4 Standing wave displayed in the bottom graph is the sum of the two
counter-propagating traveling waves displayed in the top two graphs.
207 Pulses on strings
t
into the wave equation. First we take the partial derivatives:
∂y(x, t) ∂2 y(x, t)
= −3 f 0 (x − 3t) + 3g0 (x + 3t) , = 32 f 00 (x − 3t) + 32 g00 (x + 3t) (7.54)
∂t ∂t2
∂y(x, t) ∂2 y(x, t)
= f 0 (x − 3t) + g0 (x + 3t) , = f 00 (x − 3t) + g00 (x + 3t) . (7.55)
∂x ∂x2
Here we have assumed that the derivative f 0 , f 00 , g0 , g00 exist, which means that f , f 0 , g,
and g0 must be continuous functions. Dividing Eq. (7.54)b through by 32 , we see that it is
equal to Eq. (7.55)b. That is,
1 ∂2 y(x, t) ∂2 y(x, t)
= , (7.56)
32 ∂t2 ∂x2
which is equivalent to the wave equation Eq. (7.8).
7.4 Pulses on strings
We just showed that Eq. (7.53) represents a general solution to the wave equation where
the functions f and g are any continuous functions with continuous first derivatives. We
can use this result to construct pulse solutions that propagate on a string. For example, the
following functions, plotted in Fig. 7.5, are all solutions3 to the wave equation provided
ξ = x ± 3t:
2
/2a2
fa (ξ) = Ae−ξ (7.57a)
2 2
fb (ξ) = Be−ξ /2b cos kξ (7.57b)
−2[ln(1−2ξ/d)]2 1

De
 if ξ/d < 2
fc (ξ) =  (7.57c)

1
0
 if ξ/d ≥ 2
−2[ln(1+2ξ/g)]2 1

Ge
 if ξ/d > 2
fd (ξ) =  . (7.57d)

1
0
 if ξ/d ≤ 2
For fa (ξ) and fc (ξ), we set ξ = x − 3t, and for fc (ξ) and fd (ξ), we set ξ = x + 3t in Fig.
7.5 so that the pulses propagate to the right or left, respectively, as indicated by the arrows.
You cannot tell which direction a pulse is propagating from the shape of the pulse; that is
determined solely by the sign in ξ = x ± 3t.
On the other hand, you can say something about the force on a pulse simply by examin-
ing the shape of the pulse. The force F per unit length on a segment of string is given by
mass per unit length times its (vertical) acceleration, which according to Eq. (7.9), is
∂2 y ∂2 y
F =µ = T . (7.58)
∂t2 ∂x2
Thus, the force per unit length is proportional to the curvature ∂2 y/∂x2 of the pulse. Figure
3 You can check that all these functions are continuous with continuous first derivatives.
208 Strings
t
(a) (b)
x x
−2 −1 1 2 −2 −1 1 2
(c) (d)
x x
−2 −1 1 2 −2 −1 1 2
t
Fig. 7.5 Traveling pulses.
7.6 shows a Gaussian pulse given by Eq. (7.57a) together with vectors indicating the force
at various positions along the pulse. The vertical dotted lines indicate the points where the
curvature changes sign. At the forward edge of the pulse where x > 0.25, the force is in the
same direction as the vertical velocity of the string, which serves to accelerate the vertical
displacement of the string causing the string to rise as the peak in the pulse approaches.
The restoring force is largest where the curvature is largest, which also corresponds to the
point where the displacement of the string is greatest: x = 0, at this particular moment
in time. For x < 0.25, the force is positive but the velocity is negative. This means that
the force is acting to decelerate that section of the string, which will soon come to rest at
y = 0 as the pulse propagates on down the string. This same scenario continues as the pulse
moves to the right along the string.
7.5 Energy of a vibrating string
When a pulse propagates along a string, it carries with it kinetic and potential energy. The
kinetic energy is associated with the transverse velocity of a string that has mass µ per unit
length. The potential energy is associated with the energy stored in the string when it is
stretched from its equilibrium position.
One of the fascinating properties of waves is that they can transport energy from over
large distances even though the medium supporting the wave—in this case a string—never
moves very far at all. Consider, for example, the Gaussian pulse propagating in the +x
direction depicted in Fig. 7.6. As the pulse—the disturbance—propagates along the string,
energy is transported with the disturbance, which can travel along the full length of the
string. Individual elements of the string, however, only move up and down, and never
209 Energy of a vibrating string
t
y
x/ a
−3 −2 −1 1 2 3
t
Fig. 7.6 Traveling Gaussian pulse shown as solid line. Dashed line shows pulse a moment
after the peak passes x = 0. Arrows show forces on string. Filled gray curve shows
energy density along string.
move very far from their equilibrium positions. This behavior is to be contrasted with the
energy transported by a particle—say a projectile—where the kinetic and potential energy
are transported with the particle. Pulses on strings, however, transport energy without trans-
porting material. This is a very general property of all kinds on waves: waves on strings,
sound, surface waves on water, and electromagnetic waves, to name just a few.
An interesting property of waves on strings is that the disturbance associated with a
wave, the displacement of individual mass elements, is perpendicular to the direction in
which energy is transported. Such a wave is said to be transverse. Surface waves on water
and electromagnetic waves are also transverse waves. Sound waves in air, by contrast, are
longitudinal waves: the disturbance back and forth of the molecules that carry the sound is
along the same axis as the propagation of energy.
7.5.1 Energy in a pulse
To calculate the kinetic energy, we consider mass segments of length dx. Each segment has
mass µdx and velocity ∂y/∂t. The kinetic energy of a segment of string is 21 µdx (dy/dt)2 ,
and the total energy kinetic energy for a string of length L is
L !2
1 ∂y
Z
K= µ dx . (7.59)
2 0 ∂t
Potential energy is stored when a segment of length dx is stretched to a length ds, as

shown in Fig. 7.7. We can determine this potential energy by calculating the work required
to stretch such a segment. The work is done by the constraints holding the string at a
tension T . Since the tension always acts along the tangent of the string, the work done on
a segment of string is T (ds − dx). The projected length of the string segment along x is dx
210 Strings
t
and along y is dy = (∂y/∂x) dx. Therefore, the length of the segment is
 !2 1/2  !2 

2 2 1/2
∂y  1 ∂y 
ds = dx + dy = dx 1 +  ≈ dx 1 +  ,
 
(7.60)
∂x  2 ∂x 
 
where we have neglected higher order terms in the expansion of the square root because
the string displacements are small so that the slope ∂y/∂x 1. The total work done, and
thus the total potential energy of the stretched string, is obtained by integrating over the
length of the string
Z L !2
1 ∂y
Z
U=W= T (ds − dx) = T dx . (7.61)
2 0 ∂x
The total energy is the sum of the kinetic and potential energies
!2 !2 
1 L  ∂y
Z  Z L
∂y 
E= µ +T  dx = u(x, t) dx , (7.62)
2 0  ∂t ∂x  0
where u(x, t) is the energy density, the energy per unit length of the string at a point x and
time t, which is defined as
 !2 !2   !2 !2 
1  ∂y ∂y  1  1 ∂y ∂y 
u(x, t) ≡ µ +T = T +  . (7.63)
2 ∂t ∂x  2  32 ∂t ∂x 
7.5.2 Energy density of a traveling pulse
Let’s examine the energy density of a traveling wave propagating in the +x direction,
which, as shown in §7.3, can be written as f (x − 3t). Substituting f (x − 3t) into Eq. (7.63)
and noting that
∂y ∂y
= −3 f 0 (x − 3t) , = f 0 (x − 3t) , (7.64)
∂t ∂x
we find
uR (x, t) = T f 0 (x − 3t) 2 , (7.65)

y
0.02
ds 0.01
dy
dx
x
t
− 1.0 − 0.5 0. 0 0.5
Fig. 7.7 Geometry for calculating potential energy of a string segment. Note that the y scale is
greatly magnified compared to the x scale.
t
where uR (x, t) indicates the energy density of a wave traveling to the right. Obviously, we
can perform the same analysis for waves propagating in the −x direction with the result
uL (x, t) = T f 0 (x + 3t) 2 . (7.66)

In what follows, we illustrate some properties of pulses using right-going pulses but the
same ideas apply equally to left-going waves.
As our first example, consider the Gaussian pulse given by Eq. (7.57a). Noting that
2
/2a2
f 0 (ξ) = −(ξ/a2 ) Ae−ξ (7.67)
we obtain the energy density of the pulse from Eq. (7.65)

(x − 3t)2 2 −(x−3t)2 /a2
uR (x, t) = T Ae . (7.68)
a4
The gray filled curve in Fig. 7.6 shows the energy density along the string. Notice that it is
largest where the absolute value of the slope of the string is greatest, as expected from Eq.
(7.65).
To obtain the total energy of the pulse, we integrate over all x. Since the total energy in
the pulse does not change with time, we can choose any convenient time to perform the
integration. Thus, setting t = 0 and integrating over all x, we obtain
√
x2 2 −x2 /a2 π A2
Z ∞
Epulse = T 4A e dx = T . (7.69)
−∞ a 2 a
The energy of the pulse is proportional the the square of the amplitude A2 , which is the
same familiar amplitude dependence of the energy we noted in harmonic oscillators. See,
for example, Eq. (1.52) and the accompanying discussion.
The energy of the pulse is also inversely proportional to a, the width of the pulse. This
means that it takes more energy to create a short pulse than a longer pulse with the same
amplitude. This a general property of waves on strings and all other kinds of waves. It’s
an important fact to keep in mind: confining a wave to a smaller region of space generally
requires and stores more energy.
Applying the same analysis to the sinusoidal Gaussian pulse given by Eq. (7.57b) and
depicted in Fig. 7.5b, yields the following expression for the total energy of the pulse:
√ √ 
π A2 h 1 −k2 b2

2 2
i π 2 1/b if kb 1

Epulse = 1 + e + k b A k2 b if kb 1 . (7.70)

T 2 → T
2 b 2


Here we see that the total energy stored in the pulse depends not only on the pulse width b
but also on the wave vector k. When the wavelength 2π/k is comparable to or smaller than
the pulse width b, or kb & 1, fb (x − 3t) exhibits oscillations within the Gaussian envelope.
These oscillations require energy to create (and thus also store energy). This extra energy
is reflected in the k2 b2 terms in Eq. (7.70), especially the latter one that grows quadratically
as k increases. On the other hand, when the wavelength is much longer than the width of
the Gaussian envelope, or kb 1, fb (x−3t) approaches the same form as fa (x−3t). Indeed,
in the limit that kb 1, Eq. (7.70) reduces to Eq. (7.69) when the pulse widths of the two
Gaussians are set equal, that is, when a = b.
212 Strings
t
The most important message to take away from this analysis is that creating more undu-
lations in a waveform costs (and stores) energy.
7.5.3 Transport of energy
One of the most important physical consequences of pulse propagation is that energy is
transported from one part of the string to another. To examine this, we consider the time
rate of change of energy at a point along the string as given by Eq. (7.63)
  !2 !2 
∂u ∂  1  1 ∂y ∂y 
=  + (7.71)
 
T  2

∂t ∂t 2 3 ∂t ∂x 
 
 
 !2 !2 
1  1 ∂ ∂y ∂ ∂y 
= T  2 + (7.72)
2 3 ∂t ∂t ∂t ∂x 

1 1 ∂y ∂2 y ∂y ∂ ∂y
" ! ! #
= T 22 + 2 (7.73)
2 3 ∂t ∂t2 ∂x ∂t ∂x
∂y ∂2 y ∂y ∂ ∂y
" ! ! #
=T + (7.74)
∂t ∂x2 ∂x ∂t ∂x
∂ ∂y ∂y ∂y ∂ ∂y ∂y ∂ ∂y
" ! ! # !
=T − + (7.75)
∂x ∂t ∂x ∂x ∂x ∂t ∂x ∂t ∂x
∂ ∂y ∂y
!
=T , (7.76)
∂x ∂t ∂x
where we have used the wave equation Eq. (7.8) in going from Eq. (7.73) to Eq. (7.74),
and we have used the identity (obtained using the chain rule)
∂ ∂y ∂y ∂y ∂2 y ∂y ∂ ∂y
! ! !
= + (7.77)
∂x ∂t ∂x ∂t ∂x2 ∂x ∂x ∂t
in going from Eq. (7.74) to Eq. (7.75). Defining
∂y ∂y
!
S x = −T , (7.78)
∂t ∂x
we can write Eq. (7.76) succinctly as
∂u ∂S x
=− (7.79)
∂t ∂x
For waves moving in the +x direction, y = f (x − 3t) so that
∂y ∂y
!
S x = −T = −T (−3 f 0 ) f 0 = 3uR , (7.80)
∂t ∂x
where the last equality follows from Eq. (7.65). Equation (7.80) suggests a physical in-
terpretation for S x , namely that it is the flux or current 3u of energy flowing in the +x
direction. The energy current in the −x direction would be 3uL where with uL is given by
Eq. (7.66). Thus, if there are both right and left propagating waves, the net energy flux to
t
the right would be
∂y ∂y
!
S x = −T = 3 (uR − uL ) = S R + S L . (7.81)
∂t ∂x
To better understand the physical meaning of Eq. (7.79), consider the rate of change of
energy over a segment of string between x1 and x2 :
Z x2 Z x2
∂u ∂S x
! !
dx = − dx = S (x1 ) − S (x2 ) . (7.82)
x1 ∂t x1 ∂x
This equation says that the rate of change of the energy in the line segment from x1 to
x2 is equal to the energy flux S x (x1 ) entering the segment at x1 minus the energy flux
S x (x2 ) leaving at x2 . Thus, Eq. (7.82), and by extension Eq. (7.79), are statements of energy
conservation for waves propagating on a string.
To illustrate these ideas, let’s return to our example of a Gaussian pulse, given by Eq.
(7.57a), propagating to the right along a string. We will test Eq. (7.82) over the interval
from x = 0 to x = a at t = 0. We start from the expression for u(x, t) given by Eq. (7.68).
In order to simplify calculating the derivative needed to evaluate the left hand side of Eq.
(7.82), we write it as a function of ξ = x − 3t
ξ2 2 −ξ2 /a2
uR (ξ) = T Ae . (7.83)
a4
Then, we have
∂u(x, t) ∂ξ ∂u(x, t) A2 2x x2 −x2 /a2
! ! ! !
= = −3 T 1 − e . (7.84)
∂t t=0 ∂t t=0 ∂ξ t=0 a2 a2 a2
With this result, the left hand side of Eq. (7.82) becomes
Z a
∂u(x, t) A2 a 2x x2 −x2 /a2
! Z !
dx = −3 T 2 1− 2 e dx (7.85)
0 ∂t t=0 a 0 a2 a
A2 1
Z
y = x2 /a2
h i
= −3 T 2 (1 − y) e−y dy , (7.86)
a 0
1
A2 A2
= −3 T 2 y e−y = −3 T 2 e−1 . (7.87)

a 0 a
Starting again from Eq. (7.68), the right hand side of Eq. (7.82) is
x2 2 −x2 /a2

S x (x1 , t) = 3 u (0, 0) = 3 T 4 A e = 0 (7.88)
a
2
x=0
x 2 2 A2
S x (x2 , t) = 3 u (a, 0) = 3 T 4 A2 e−x /a = 3 T 2 e−1 , (7.89)
a x=a a
which means that
A2 −1
S x (x1 , t) − S x (x2 , t) = −3 T
e . (7.90)
a2
The equality of Eqs. (7.87) and (7.90) demonstrates that the change in energy along a seg-
ment of string is equal to the next flux of energy into that segment, as state mathematically
214 Strings
t
by Eq. (7.82). This can be understood graphically by considering how the gray area, which
represents the energy density, between 0 and a in Fig. 7.6 changes as the pulse pulse moves
infinitesimally to the right. In this case, no energy moves into the region at x = 0 as the
energy density u there is zero. By contrast, at x = a, u is finite and thus area is lost as the
pulse propagates to the right. Thus, at t = 0, energy is leaving the interval x = [0, 1] as the
pulse passes by.
7.5.4 Energy density of a traveling wave
We consider now a sinusoidal traveling wave of the form given by Eq. (7.40) where y(x, t) =
A sin (kx − ωt − ϕ) and calculate the energy density using Eq. (7.63)
 !2 !2 
1  1 ∂y ∂y 
u(x, t) = T  2 + (7.91)
2 3 ∂t ∂x 

1 1
" #
= T 2 −ωA cos (kx − ωt − ϕ) 2 + kA cos (kx − ωt − ϕ) 2 (7.92)

2 3
= T k2 A2 cos2 (kx − ωt − ϕ) (7.93)
We obtain the average energy by averaging over one entire period T
1 T 1 T 1
Z Z
ū = u(x, t) dt = T k2 A2 cos2 (kx − ωt − ϕ) dt = T k2 A2 . (7.94)
T 0 T 0 2
Once again we see that the energy is proportional to A2 . It is also proportional to k2 , which
is the same dependence we obtained for pulses in Eq. (7.70) in the limit of large k. Large
values of k correspond to short wavelengths and more undulations, which as stated in the
previous section, involves more energy.
The average energy flux for a traveling wave of the form y(x, t) = A sin (kx − ωt − ϕ) is
give by
1
S x = 3 ū = 3 T k2 A2 (7.95)
2
7.5.5 Energy of standing waves
Standing waves can be represented as a sum over normal modes

∞
X
y(x, t) = An cos(ωn t − ϕn ) sin kn x . (7.96)
n=1
where ωn = kn 3 and kn = πn/L for fixed boundary conditions or kn = 2πn/L for periodic
boundary conditions. The amplitude of mode n is An .
In §7.2.4 we learned that standing waves consist of counter-propagating sinusoidal waves.
Those counter-propagating waves carry equal fluxes of energy, but in opposite directions.
Therefore, the net energy flux to the right or the left of standing waves is zero.
Because the net energy flux is zero, we will simply calculate the total energy of standing
215 Boundaries
t
waves by integrating Eq. (7.63) over the entire length of the string from x = 0 to L. The
partial derivatives appearing in Eq. (7.63) are given by
∞
∂y(x, t) X
=− ωn An sin(ωn t − ϕn ) sin kn x (7.97)
∂t n=1
∞
∂y(x, t) X
= kn An cos(ωn t − ϕn ) cos kn x (7.98)
∂x n=1
To calculate the energy we need to integrate (∂y/∂t)2 and (∂y/∂x)2 over x from 0 to L.
These squares are double sums over indices n and m of the form
Z L !2 Z L
∂y X
dx = S n (t)S m (t) sin kn x sin km x dx (7.99a)
0 ∂t n,m 0
Z L !2 Z L
∂y X
dx = Cn (t)Cm (t) cos kn x cos km x dx , (7.99b)
0 ∂x n,m 0
where S n (t) = ωn An sin(ωn t − ϕn ) and Cn (t) = kn An cos(ωn t − ϕn ). The integrals can be

evaluated with the help of the orthogonality relations
Z L Z L
2
sin kn x sin km x dx = cos kn x cos km x dx = δn,m . (7.100)
0 0 L
where δn,m is the Kronecker-δ, which is equal to 1 if n = m and zero otherwise. Thus, the
double sums in Eq. (7.99) reduce to a single sum over n. Then, because µω2n = µ32 kn2 =
T kn2 , the total energy reduces to
Z L  !2 !2 
1  1 ∂y ∂y 
E= T +  dx (7.101)
3 ∂t ∂x 
 2
2 0
∞
1 L X 2 2h 2
kn An sin (ωn t − ϕn ) + cos2 (ωn t − ϕn )
i
= T (7.102)
2 2 n=1
∞ ∞
L X 2 2 L X 2 2
= T kn An = µ ω A (7.103)
4 n=1 4 n=1 n n
Note that the contribution of each mode is independent of all the others. The energy is the
sum of the contributions from each mode. If the amplitude of a given mode is zero, it does
not contribute to the energy.
7.6 Boundaries
216 Strings
t
Problems
7.1
8 Sound and Electromagnetic Waves
8.1 Sound
8.2 Reflection and refraction
8.3 Electromagnetic waves
217
218 Sound and Electromagnetic Waves
t
Problems
8.1
Interference, Diffraction, and Fourier
9
Optics
9.1 Interference
9.2 Diffraction
9.3 Fourier optics
9.4 Scattering
219
220 Interference, Diffraction, and Fourier Optics
t
Problems
9.1
PART II
CLOSING FEATURES
A Appendix A Linear algebra
In this Appendix, we review the properties of matrices that are relevant to our study of
coupled oscillators. For our purposes, matrices can be regarded as generalizations of vec-
tors, so before summarizing the properties of matrices, we briefly review vectors. We then
provide a concise overview of those aspects of linear algebra and matrices that are relevant
to our study of oscillations and normal modes.
A.1 Vectors
A vector is an ordered list of numbers. The numbers that make up a vector are called its
components. The number of components in the list is called the dimension of the vector. For
example, we can represent the spatial location of an point in space by a three-dimensional
vector, which we might express in terms of its Cartesian components: r = (r x , ry , rz ) =
(x, y, z). In general, higher dimensional vectors can be constructed, say with N components:
S = (s1 , s2 , s3 , ..., sN ). Vectors also obey a number of familiar rules pertaining to how
they are added, subtracted, and multiplied. Here we summarize those rules for the vectors
a = (a1 , a2 , a3 ) and b = (b1 , b2 , b3 ):
ta = (ta1 , ta2 , ta3 ) scalar multiplication (A.1)

a · b = (a1 b1 , a2 b2 , a3 b3 ) dot product (A.2)
a × b = (a2 b3 − a3 b2 , a3 b1 − a1 b3 , a1 b2 − a2 b1 ) cross product (A.3)
Scalar multiplication and the dot product are commutative. That is, the order of multipli-
cation doesn’t matter so that ta = at and a · b = b · a. By contrast, the cross product does
not commute; in general a × b , b × a, as can be readily demonstrated using the above
definition.
Each of the three types of vector multiplication have a geometrical interpretation. Scalar
multiplication of a by t means that the vector a is stretched by a factor or t. The dot product
a · b represents a scalar with a length equal to the projection of a onto b or of b onto a. The
cross product a × b gives a vector representing the rotation of a into b.
223
224 Linear algebra
t
A.2 Matrices
Matrices are arrays of numbers arranged in rows and columns. A 2 row by 3 column matrix,
or 2 × 3 matrix, A is defined as
A11 A12 A13

!
A= . (A.4)
A21 A22 A23
We denote matrices by bold sans-serif characters and their elements by the same sub-
scripted italic characters. We shall often work with square matrices, meaning matrices that
have the same number of rows and columns
A11 A12 A13 

 
A = A21 A22 A23  . (A.5)

 
A31 A32 A33
 
The transpose of a matrix A , denoted as A T , is obtained by reversing the columns and

rows, i.e. ATij = A ji . Thus, the transpose of a matrix with 2 rows and 3 columns would have
3 rows and 2 columns.
With these definitions, we see that vectors are simply matrices that have either one row or
one column. When working with matrices, we generally write a column vector (1 column,
n rows) as a and a row vector (n columns, 1 row) as the transpose of a column vector a T .
A.2.1 Matrix operations
Matrices can be added and subtracted. Like vectors, they can also be multiplied in several
different ways.
Matrix addition
Matrix addition and subtraction are defined on an element-by-element basis. Therefore, to
be added or subtracted, two matrices must have the same number or rows and columns. For
the 2 × 3 matrices A and B , this means that
A11 A12 A13 B B12 B13

! !
C = A +B = + 11 (A.6)
A21 A22 A23 B21 B22 B23
A11 + B11 A12 + B12 A13 + B13
!
= , (A.7)
A21 + B21 A22 + B22 A23 + B23
which we can rewrite more compactly in component form as Ci j = Ai j + Bi j . Matrix

subtraction for C = A − B is defined in the obvious way Ci j = Ai j − Bi j .
225 Properties of matrices
t
Matrix multiplication
As in the case of vectors, several types of multiplication are defined for matrices. In scalar
multiplication, all components of a matrix are multiplied by the same scalar value: the
components of tA , where t is a scalar, are given by t Ai j for all indices i j.
The matrix dot product C = A · B , often simply called matrix multiplication and written
C = AB , is defined by the rule
X
Ci j = Aik Bk j , (A.8)
k
or, using the summation convention, simply as

Ci j = Aik Bk j , (A.9)
where it is understood that repeated indices (in this case k) are summed over. This definition
only makes sense, and thus the matrix dot product is only defined, if the number of columns
in A is equal to the number of rows in B . The matrix dot product is also often referred to
as the inner product. Except in special cases, matrix multiplication is not commutative,
that is, AB , BA in general. Matrices that do commute must be square, because of the
definition of matrix multiplication. This is a necessary but not sufficient condition for two
matrices to commute.
The matrix outer or direct product is also defined, and is usually designated as A ⊗ B .
But as we will not need it, we do not discuss it further.
A.3 Properties of matrices
Just as vectors have a number of properties, such as length and direction, matrices also
have a number of properties that can be characterized by various well-defined quantities.
Here we review a few of the more relevant of those properties.
A.3.1 Determinant of a matrix
The determinant is defined for square matrices and is designated by absolute value signs
(or in some texts as det A ):
A11 A12 A13 A1n

···
A
21 A22 A23 A2n

···
|A | = det A = A31 A32 A33 ··· A3n (A.10)
. .. .. .. ..
.. . . . .
An1 An2 An3 Ann

···
The determinant of a 2 × 2 matrix is given by

A11 A12
|A | = = A11 A22 − A12 A12 (A.11)
A12 A22
226 Linear algebra
t
The determinant of any higher order matrix can be calculated using Laplace’s formula
n
X
|A| = (−1)1+ j A1 j M1 j , (A.12)
j=1
where M 1 j is the (n−1)×(n−1) matrix formed by removing the first row and the jth column
of A , the so-called 1 j minor of A . Thus for a 3 × 3 matrix, the determinant is given by
A11 A12 A13

|A | = A21 A22 A23 (A.13)

A31 A32 A33

A A23 A21 A23 A21 A22
= A11 22 − A12 + A 13 (A.14)
A32 A33 A31 A33 A31 A32
= A11 (A22 A33 − A23 A32 ) − A12 (A21 A33 − A23 A31 + A13 (A21 A32 − A22 A31 ) .(A.15)
Larger matrices can be broken down in terms of their minors in a similar fashion. In fact,
there is no need to use first row as the special row in the definition of |A | as we did in
Eq. (A.12). Any row or column may be designated as the “special" row or column around
which the determinant is calculated, but it is often convenient to choose the first row, as we
have done above.
Exercise A.3.1 Using the definition of the determinant given by Eq. (A.12), show that
the the determinant of a lower or upper triangular matrix is given by the produce
of its diagonal elements. That is, show that |L | = ni=1 Lii and that |U | = ni=1 Uii ,
P P
where L is a lower triangular matrix, meaning that all the elements above the matrix
diagonal are zero, and where U is a upper triangular matrix, meaning that all the
elements below the matrix diagonal are zero.
A.3.2 Inverse of a matrix
The inverse of a matrix A is designated as A −1 and has the property that A −1A = I where I
is the identity matrix.
A.3.3 Matrix symmetry
A matrix A is said to be symmetric if its elements Ai j obey the relation Ai j = A ji . Thus, for a
symmetric matrix A T = A . Similarly, a matrix B is said to be antisymmetric if its elements
Bi j obey the relation Bi j = −B ji . Thus, for an antisymmetric matrix B = −B T . Clearly, for
these relations to hold, the matrices A and B must each have the same numbers of rows and
columns, so only square matrices can be symmetric or antisymmetric.
Of course, most square matrices are neither symmetric or antisymmetric. However, it’s
easy to show that any square matrix can be expressed as the unique sum of a symmetric
and antisymmetric matrix. See if you can show that.
227 Properties of matrices
t
A.3.4 Hermitian matrices
A matrix A is said to be Hermitian1 if the complex conjugate of the transpose of the matrix
is equal to the matrix itself, i.e. if A ∗ T = A . The complex conjugate of the transpose of
the matrix A is called the adjoint and is often denoted as A † , i.e. A † ≡ A ∗ T . Therefore,
Hermitian matrices are matrices for which A † = A and are thus said to be self-adjoint.
A.3.5 Positive definite matrices
An n × n square matrix P is said to be positive definite if a T P a > 0 for all vectors a with n
real components. If a contains complex components, then P is said to be positive definite
if a†P a > 0 for all vectors a with n complex components, where a† is the transpose of
the complex conjugate of a. Clearly, the first definition is a special case of the second. It
is straightforward to show that all Hermition matrices are positive definite. Because a real
symmetric matrix is Hermitian, real symmetric matrices are also positive definite.
Positive definite matrices have a number of useful properties. Below we list those that
are particularly useful for our purposes.
• The eigenvalues of positive definite matrices are all real and positive.
• Positive definite matrices P can always be factorized such that P = LL † where L is a
unique lower triangular matrix with diagonal elements that are all positive. This factor-
ization is called Cholesky decomposition. If P is real, then P = LL T where L is real.
• Two positive definite matrices can be simultaneously diagonalized, that is, diagonalzed
by the same set of similarity transformations.
A.3.6 Orthogonal matrices
A matrix A is said to be orthogonal if its inverse is also its transpose, i.e. if A −1 = A T . Of

course this means that for an orthogonal matrix AA T = I . The determinant of an orthogonal
matrix is either 1 or -1; that is, |A| = ±1.
A.3.7 Matrix identities
We list here a few useful matrix identities, most of which are easily proved:
(AB )−1 = B −1A −1 (AB )T = B T A T (A −1 )T = (A T )−1

|AB | = |A | |B | |A −1 | = 1/|A |
1 after the 19th century French mathematician Charles Hermite

228 Linear algebra
t
A.4 Transformations
Matrices that operate on (e.g. multiply) vectors or other matrices represent linear transfor-
mations. The kinds of transformations represented by matrix operations can be classified
into a number of different types, many of which have important geometrical or mathemat-
ical properties. Below we consider a few that are useful for our studies.
A.4.1 Rotations
A.4.2 Choleski decomposition
In linear algebra it is often useful to write a square matrix A as the product of a lower
triangular matrix L and an upper triangular matrix U . That is,
A = LU (A.16)
(A.17)
or
A11 A12 A1n  L11 0 0  U11 U12 U1n 
    
··· ··· ···
A A22 ··· A2n  L21
 
L22 ··· 0   0 U22 ··· U2n 

 21
 . .. ..  =  .. .. ..   . .. ..  .
 
.. .. ..
  
 .. . . .   . . . .   .. . . . 
An1 An2 Ann Ln1 Ln2 Lnn 0 0 Unn
  
··· ··· ···
  
Factoring a matrix this way is known as LU decomposition and is generally useful in nu-
merical methods for solving systems of linear equations. Such a factorization is not unique;
there are many LU decompositions for a given matrix. For square symmetric matrices S
that are positive definite (§A.3.5), it can shown that the LU decomposition can be done
such that U = L T , that is
S = LL T . (A.18)
This special case of LU decomposition where S = LL T is called Choleski decomposition.

Let’s start by explicitly calculating L for a general 3 × 3 matrix:
S = LL T (A.19)
or
S 11 S 21 S 31  L11 0 0  L11 L21 L31 
    
S 21 S 22 S 32  = L21 L22 0   0 L22 L32  (A.20)
     
S 31 S 32 S 33 L31 L32 L33 0 0 L33
    
 2
 L11 L11 L21 L11 L31


2 2
= L11 L21 L21 + L22 L21 L31 + L22 L32  . (A.21)
 
2 2 2 
L11 L31 L21 L31 + L22 L32 L31 + L32 + L33

229 Transformations
t
Note that the matrix LL T is symmetric, as it should be. We can determine the elements of
L by equating S element by element with LL T starting with the first column. Thus we have
2
p
S 11 = L11 ⇒ L11 = S 11 (A.22)
S 21 = L11 L21 ⇒ L21 = S 21 /L11 (A.23)
S 31 = L11 L31 ⇒ L31 = S 31 /L11 (A.24)
Notice that Eqs. (A.23) and (A.24) are obtained using L11 the result of Eq. (A.22). Contin-
uing with the second column, starting with the diagonal element we obtain
q
2 2 2
S 22 = L21 + L22 ⇒ L22 = S 22 − L21 (A.25)
S 32 = L21 L31 + L22 L32 ⇒ L32 = (S 32 − L21 L31 )/L22 . (A.26)
Once again we see that Eqs. (A.25) and (A.26) both use results from previous calculations
of various matrix elements of L . Finally, from the third column we obtain
q
2 2 2 2 2
S 33 = L31 + L32 + L33 ⇒ L33 = S 33 − L31 − L32 . (A.27)
Generalizing these results to an N × N matrix, the diagonal terms are given by
v
j−1
u
t
X
Ljj = S jj − L2jk . (A.28)
k=1
Similarly, the off-diagonal terms are given by

j−1
 
1  X 
Li j = S i j − Lik L jk  . (A.29)
Ljj k=1
Equations (A.28) and (A.29) are suitable for using in a computer routine to numerically
compute the matrix L for any given S . Our assumption that S is a positive definite matrix
guarantees that the argument of the square root in Eq. (A.28) is always positive. Here is a
Python function that implements Eqs. (A.28) and (A.29) to calculate L :
import numpy as np
def choleski(s):
n = len(s)
for k in range(n):
try:
s[k,k] = np.sqrt(s[k,k] - np.dot(s[k,0:k],s[k,0:k]))
except ValueError: # for square root of negative number
print("Matrix is not positive definite")
for i in range(k+1,n):
s[i,k] = (s[i,k] - np.dot(s[i,0:k],s[k,0:k]))/s[k,k]
for k in range(1,n): # put zeros in upper triangle
s[0:k,k] = 0.0
return s
230 Linear algebra
t
The program performs the Choleski decomposition “in place" meaning that it puts the
result into the same matrix that was input into the routine. This has the virtue of being
faster and saving computer memory. However, if you want to save the original matrix,
you had better make a copy of it using the NumPy copy function before performing the
Choleski decomposition.
Exercise A.4.1 Perform Choleski decomposition of the following symmetric matrix
by hand using the algorithm described by Eqs. (A.28) and (A.29). Verify that the
algorithm works in this case by explicitly showing that S = LL T .
 6 20 15
 
S = 20 80 50
 
15 50 60
 
Exercise A.4.2 Use the Python function given above to find the Choleski decompo-
sition of the matrix S above. Verify that you get the same answer that you got in
Problem 1 above. Explain how the program works, line by line. Include in your ex-
planation why it is not a problem that this routine calculates L in place.
A.4.3 Similarity transformations
A.5 Eigenvalue equations
The solution to a system of homogeneous linear equations of the form C a = 0 exists only
if det C = 0. We can easily prove this for a general 2 × 2 matrix C with matrix elements Ci j
and column vector a with elements a1 and a2 by writing out the matrix equation C a = 0
explicitly
C11 C12 a1
! !
=0. (A.30)
C21 C22 a2
or
C11 a1 + C12 a2 = 0 (A.31)
C21 a1 + C22 a2 = 0 (A.32)
From Eq. (A.32),
C21
a2 = − a1 . (A.33)
C22
Substituting this into Eq. (A.31) yields
C21
C11 a1 − C12 a1 = 0 (A.34)
C22
(C11C22 − C12C21 ) a1 = 0 (A.35)
det C ≡ (C11C22 − C12C21 ) = 0 . (A.36)
231 Eigenvalue equations
t
The last step assumes that a1 , 0. If a1 = 0, then a2 = 0 provided C21 , 0. This is an
uninteresting case in which a1 = a2 = 0. If C21 = 0, the C22 a2 = 0 and either C22 = 0 or
a2 = 0. The latter case is uninteresting again with a1 = a2 = 0. If C22 = 0 and C21 = 0, then
det C = 0, which is uninteresting. The end result is that either a1 = a2 = 0 or det C = 0.
Lagrangian formulation of classical
B
mechanics
You are used to finding the equation of motion in classical mechanics by careful application
Newton’s second law, F = ma. But “there is more than one way to skin a cat" as the
saying goes. In this appendix, we show you a versatile and powerful method for finding
the equation of motion for a dynamical system.
The method was developed in the 18th century by Joseph Louis Lagrange and constitutes
an elegant and powerful reformulation of Newtonian mechanics. The method is particularly
useful when there are elements in the problem that constrain the motion of the system, as
the method eliminates the need to consider the forces associated with those constraints. An
example is the rod in a pendulum that constrains the mass at the end of the rod to move in
an arc. The Lagrangian formulation eliminates the need to consider the tension in the rod: it
simply never appears in the formulation or solution of the problem. This approach becomes
particularly useful in more complex problems, such as the double pendulum presented
below, in which the constraint forces can be quite a nuisance to determine.
B.1 Lagrangian
The Lagrangian formulation of classical mechanics begins with a quantity called the La-
grangian L, which is the difference the kinetic energy K and potential energy U of a system.
L= K−U . (B.1)
For a simple unconstrained particle, the kinetic energy in Cartesian coordinates is given by
1 2
m ẋ + ẏ2 + ż2 .

K= (B.2)
2
As the potential energy is a function of the coordinates x, y, and z, the Lagrangian is thus a
function of the coordinates and velocities of the particle: L = L(x, y, z, ẋ, ẏ, ż). The positions
(x, y, z) and velocities ( ẋ, ẏ, ż) are the independent variables of the Lagrangian. The force
on the particle is given by the gradient of L, or in component form
∂L ∂U
=− = Fx , (B.3)
∂x ∂x
232
233 Lagrangian
t
with similar equations for Fy and Fz . Taking derivatives of L with respect to the velocities
gives the various components of the momentum
∂L
= m ẋ = p x , (B.4)
∂ ẋ
with similar equations for py and pz . Noting that Newton’s second law is most generally
given by F = dp/dt we can use these results to rewrite Newton’s second lawn terms of the
Lagrangian function L as
∂L d ∂L ∂L d ∂L ∂L d ∂L
= , = , = . (B.5)
∂x dt ∂ ẋ ∂y dt ∂ẏ ∂z dt ∂ż
Thus far, in writing down Eq. (B.5) we have just managed to rewrite F = ma in a somewhat
more complicated form. We have yet to see the utility or power of this formulation of
Newton’s second law. The power of this formulation only becomes apparent when we
consider constrained systems whose motion is described by generalized coordinates. Don’t
expect that last sentence to make much sense yet. Just be patient and follow the argument
laid out below.
Suppose we wish to describe the state of our system using some set of coordinates other
than Cartesian coordinates. For example, we may wish to use angular coordinates, as we
do for the pendulum in Chapters 1–3. For more complex problems, it is often convenient
to use a mixture of translational and orientational coordinates. For such cases, it would be
useful to reformulate Eq. (B.5) in terms of generalized coordinates.
As a first step, let us recognize that it will take a certain minimal number of coordinates
to specify the configuration of a system. For the simple pendulum described in Chapters
1–3, that number is 1, say the angle θ or the arc length s, as defined in Fig. 1.1. Of course,
we can also describe the pendulum’s position in terms of x and y, but we can eliminate the
y coordinate in favor x using the geometry specified by the constraint that the pendulum
move along an arc of constant radius l, the length of the pendulum. In other cases, say
for the coupled double pendulum discussed in §4.1.1, two coordinates may be required. In
other cases, still more coordinates may be required, but in general some minimum number
of generalized coordinates will be required. Let’s designate those generalized coordinates
as q1 , q2 , ..., qN . If each coordinate qi can be varied independently of all the other coordi-
nates q j , qi , then the coordinates are said to be holonomic. In this case, the number of
coordinates N is equal to the number of degrees of freedom of the system.
In some cases, the system of coordinates do not all vary independently. The classic ex-
ample is that of a ball rolling without slipping on a surface. In such a case, two coordinates
are required to specify translation of the center of mass of the sphere and three more (an-
gular coordinates) are required to specify its orientation. However, if the ball moves, at
least two coordinates, one for translation and one for orientation change. Moreover, there
is no simple way to eliminate one coordinate in favor of another in this case. Such coordi-
nates, which cannot all vary independently, are said to be nonholonomic. Systems requiring
nonholonomic coordinates are generally much more difficult to treat. Fortunately, we shall
need to consider only systems that can be described with holonomic coordinates. In what
follows, therefore, we assume that we can find coordinates that are holonomic.
With this in mind, we can specify the state of a system with N degrees of freedom in
234 Lagrangian formulation of classical mechanics
t
terms of N Cartesian coordinates {ri }, which would encompass the x, y, and z coordinates
of all the particles in the system. The kinetic energy of such a system is given by
N
X 1
K= mi ṙi2 , (B.6)
i
2
where the sum extends over all N degrees of freedom. In order to generalize Eq. (B.5),
we would like to express the kinetic energy in terms of the generalized coordinates {qk }.
In general, we can express each of the N Cartesian coordinates ri as functions of the N
generalized coordinates {qk }:
ri = ri (q1 , q2 , ..., qN , t) , (B.7)
where we have also included an explicit time dependence for the case the dependence of
the Cartesian and generalized coordinates is time dependent (in some known way). Thus,
X ∂ri ∂qk ∂ri
ṙi = + (B.8)
k
∂qk ∂t ∂t
X ∂ri ∂ri
= q̇k + (B.9)
k
∂qk ∂t
Taking the derivative ∂ṙi /∂qm picks out of the sum in Eq. (B.9) only that term where k = m,
as all other derivatives are zero, yielding
∂ṙi ∂ri
= . (B.10)
∂q̇m ∂qm
Of course, the index label is arbitrary so we can also write the above equation with m = k.
Multiplying both sides by ṙ and taking the time derivative gives
d ∂ṙi d ∂ri
! !
ṙi = ṙi (B.11)
dt ∂q̇k dt ∂qk
∂ri ∂ṙi
= r̈i + ṙi . (B.12)
∂q̇k ∂qk
As it doesn’t matter in which order we perform the differentiations, we can rewrite this as
d  ∂ ṙi2  ∂ri ∂  ṙi2 

   
 = r̈i +   . (B.13)
dt ∂q̇k 2 ∂q̇k ∂qk 2

Multiplying both sides by mi gives
d  ∂ mi ṙi2 

 2
∂r ∂  mi ṙi 
 
 = mi r̈i i +  . (B.14)
dt ∂q̇k 2 ∂q̇k ∂qk 2
 
Summing over i yields

d ∂K ∂ri ∂K
! X
= Fi + , (B.15)
dt ∂q̇k i
∂q̇k ∂qk
235 Lagrangian
t
where we have used Newton’s second law Fi = mi r̈i and Eq. (B.6) for the kinetic energy.
Rearranging terms, we can rewrite this equation as
d ∂K ∂K
!
− = Fk , (B.16)
dt ∂q̇k ∂qk
where Fk is known as a generalized force, which is defined as
X ∂ri
Fk ≡ Fi . (B.17)
i
∂qk
Equation (B.16) is completely general so long as the generalized coordinates {qk } are all
independent. A further simplification is possible if the force is conservative, meaning that
it can be written as the gradient of a potential energy, or in component form, Fi = −∂U/∂ri .
In this case the generalized force can be written as
X ∂ri X ∂U ∂ri
Fk = Fi =− (B.18)
i
∂qk i
∂ri ∂qk
Equation (B.7) tells us that we can write each Cartesian coordinate ri in terms of general-
ized coordinates, which in turn tells us that we can write the potential energy as a function
of the generalized coordinates U(q1 , q2 , ..., qN ). Taking the derivative with respect to qk
gives
∂U X ∂U ∂ri
= . (B.19)
∂qk i
∂ri ∂qk
Thus we see that the generalized force for a conservative system is
∂U
Fk = − . (B.20)
∂qk
Substituting this expression into Eq. (B.16) gives
d ∂K ∂K ∂U
!
= − . (B.21)
dt ∂q̇k ∂qk ∂qk
Finally, using our definition of the Lagrangian L = K − U, we obtain our final result
d ∂L ∂L
!
= , (B.22)
dt ∂q̇k ∂qk
where L is understood to be a function of the generalized coordinates {qi , q̇i } and we have
used the fact that U depends only on the coordinates {qi } so that ∂U/∂q̇i = 0. Equation
(B.22) expresses the Lagrangian reformulation of Newton’s second law for conservative
systems, that is, for systems without any dissipation (e.g. friction). Because it is expressed
in generalized coordinates, it can be significantly simpler to apply to a dynamical system
than a straightforward application of Newton’s second law. In the next section we illustrate
its utility with a couple of examples.
If there is dissipation in the problem, then Eq. (B.16) must be used. In fact, one can
replace K by L in Eq. (B.16) for all the conservative forces in the problem and then only
calculate the generalized force F for the nonconservative forces. We will not pursue further
the application of Lagrange’s equations to systems with nonconservative forces.
t
Before working some examples, we wish to make one final point. Just as Eq. (B.18)
defines a generalized force for a conservative system, we can also define a generalized
momentum
∂L
Pk = = mq̇k . (B.23)
∂q˙k
In the examples that follow, we see that the generalized forces and momenta turn out to be
actual forces and momenta or actual torques and angular momenta, depending on whether
the generalized coordinates associated with them are translational or orientational coordi-
nates.
B.2 Simple pendulum
Let’s see how these results work a few practical examples. We begin by considering an
almost trivial example, the undamped pendulum pictured in Fig. 1.1. There is only one
degree of freedom in this problem so we need only one generalized coordinate, which we
take to be θ. The kinetic and potential energies, expressed in terms of this generalized
coordinate, are
1 2 2
K= ml θ̇ , U = mgl(1 − cos θ) , (B.24)
2
where we have taken U = 0 to be where θ = 0. Therefore
1
L = K − U = ml2 θ̇2 − mgl(1 − cos θ) . (B.25)
2
The derivatives in Eq. (B.22) are
∂L ∂L
= ml2 θ̇ , = −mgl sin θ . (B.26)
∂θ̇ ∂θ
Thus, taking time derivative ∂L/∂θ̇ = ml2 θ̈, Eq. (B.22) gives
ml2 θ̈ = −mgl sin θ , (B.27)
which is the equation of motion for the simple pendulum (cf. Eq. (1.3)).
Notice that we obtained the equation of motion, an equation that involves forces, en-
tirely from energy considerations. No forces were used. Moreover, there was nowhere any
reference to the constraint force provided by the tension in the (massless) pendulum rod.
In this case, the generalized force is
∂U
Fθ = − = −mgl sin θ , (B.28)
∂θ
which is a torque, as one might expect since the generalized coordinate is an angle. Simi-
larly, the generalized momentum is actually the angular momentum
∂L
Pθ = = ml2 θ̇ . (B.29)
∂θ̇
237 Double pendulum
t
B.3 Double pendulum
To see the real power of the Lagrangian formulation, we turn to a more complex problem:
the double pendulum. This problem is deceptively difficult if one tries to apply Newton’s
second law in a straightforward way. As we shall see, it is quite manageable when analyzed
using the Lagrangian formulation.
The double pendulum is shown in Fig. B.1. For simplicity we assume that the motion of
the pendulum is confined to a single plane. The kinetic energy of the upper mass is 12 ml2 θ̇2 .
The kinetic energy of the lower mass is quite a bit more complicated, as its velocity is the
vector sum of the two velocities lφ̇ and lθ̇, as indicated in Fig. B.1. The resultant velocity
can be obtained from the law of cosines where the angle between the two vectors is φ − θ.
The total kinetic energy is thus given by
1 2 2 1 2 2
ml θ̇ + ml θ̇ + 2θ̇ φ̇ cos(θ − φ) + φ̇2 .

K= (B.30)
2 2
The potential energy of the system is of the upper mass is mgl(1 − cos θ) and of the lower
pendulum is mgl[(1 − cos θ) + (1 − cos φ)]. The total potential energy of the system is thus
U = mgl 2(1 − cos θ) + (1 − cos φ) . (B.31)

Therefore, the Lagrangian for the double pendulum is
L= K−U
1 2
!
2 2
= ml θ̇ + θ̇ φ̇ cos(θ − φ) + φ̇ − mgl 2(1 − cos θ) + (1 − cos φ) . (B.32)

2
t
Fig. B.1 Double pendulum.
t
It is straightforward, albeit a bit tedious, to apply Eq. (B.22), Lagrange’s equations, to ob-
tain the equations of motion for θ and φ. The result is a pair of highly nonlinear equations
that cannot be solved analytically. To obtain an analytical solution, you have to make the
usual simplifying small-angle assumption in order to linearize the equations of motion.
Knowing this in advance, it is simpler to make the small angle approximation on the La-
grangian itself. But, as Lagrange’s equations involve derivatives, it is necessary to keep
terms up to quadratic order in θ and φ (and θ̇ and φ̇). Thus, we make the approximations
cos θ ' 1 − 12 θ2 and cos φ ' 1 − 21 φ2 . We make the approximation cos(θ − φ) ' 1 as
cos(θ − φ) is multiplied by θ̇ φ̇, which is already quadratic in order, meaning that keeping
any higher order terms in the Taylor series expansion of cos(θ − φ) would lead the third and
higher order terms which would be discarded in the small angle approximation. Making
these small angle approximations, Eq. (B.32) becomes
L = ml2 θ̇2 + θ̇ φ̇ + 12 φ̇2 − mgl θ2 + 21 φ2 .

(B.33)
Lagrange’s equations of motion for this system are
d ∂L ∂L d ∂L ∂L
! !
= , = , (B.34)
dt ∂θ̇ ∂θ dt ∂φ̇ ∂φ
which become
ml2 (2θ̈ + φ̈) = −2mgl θ (B.35)
2
ml (θ̈ + φ̈) = −mgl φ . (B.36)
These equations can be rewritten in matrix form m ẍ = −k x where
2 1 2 0 θ
! ! !
2
m = ml , k = mgl , x= . (B.37)
1 1 0 1 φ
Using the trial solution x = aeiωt , the matrix equation becomes an eigenvalue equation
G a = λa where G ≡ m −1k . The inverse of m is
1 1 −1
!
m −1 = 2 (B.38)
ml −1 2
so that
g 2
!
−1
G =m k =
−1
(B.39)
l −2 2
To find the eigenvalues, we solve the secular equation
2g/l − λ

−g/l
|G − λI | = =0 (B.40)
−2g/l 2g/l − λ
!2 g 2
2g
= −λ −2 =0. (B.41)
l l
which, recalling that λ = ω2 , gives the normal frequencies
g √
r
ω= 2± 2 (B.42)
l
239 Symmetry of the mass and stiffness matrices
t
The eigenvectors are determined for each eigenvalue by the equation
(G − λI ) a = 0 . (B.43)
√ √
For λα = ω2α = (g/l)(2 − 2) and for λβ = ω2β = (g/l)(2 + 2), this gives the eigenvectors
1 1 1 1
! !
aα = √ √ , aβ = √ √ . (B.44)
3 2 3 − 2
B.4 Symmetry of the mass and stiffness matrices
The examples of coupled oscillators discussed in Chapter 4 and in §B.3 above illustrate the
general method of finding the normal modes of a system. We start by finding the linearized
equations of motion for a system, which are then cast as an eigenvalue problem. Solving the
problem then consists primarily of find the normal frequencies and eigenvectors (or normal
coordinates) of the system. The Lagrangian formulation provides a very general method
for finding the equations of motion based on first writing down the kinetic and potential
energies in terms of their generalized coordinates {qk , q̇k }. Because we are interested only
in the linearized equations of motion for the normal mode problem, we consider only small
oscillations about the equilibrium positions of all the different degrees of freedom.
For the cases of interest here, the potential energy is a function only of the general-
ized coordinates and not their time derivatives (the generalized “velocities"). Using vector
notation for the entire set of normal coordinates, q ≡ {qk }, we can express the potential
energy of the system for small excursions from equilibrium by a Taylor series about their
equilibrium positions, defined by q = 0
X ∂U ! 1 X ∂2 U
!
U(q) = U(0) + qk + qk ql + ... (B.45)
k
∂qk 0 2 k,l ∂qk ∂ql 0
Since we are expanding U about the equilibrium positions, the first derivatives (∂U/∂qk )0
are zero. Thus, to second order in the displacements, we can write the potential energy as
1X
U(q) = kkl qk ql , (B.46)
2 k,l
where
∂2 U
!
kkl ≡ , (B.47)
∂qk ∂ql 0
and we have chosen the zero of potential energy so that U(0) = 0. Note that kkl = klk , as
the order of differentiation doesn’t matter.
t
The kinetic energy is a function of the generalized velocities q̇ ≡ {q̇k }, and quite gener-
ally have the form
1X
K(q, q̇) = mkl (q) q̇k q̇l , (B.48)
2 k,l
where the coefficients mkl (q) can depend on the generalized coordinates q (see Eq. (B.30)
for an example). For small deviations from equilibrium we can expand mkl (q) in a Taylor
series about q = 0:
X ∂mkl
mkl (q) = mkl (0) + q s + ... (B.49)
s
∂q s
However, the expression for the kinetic energy Eq. (B.48) is second order in the general-
ized velocities as every term in the sum contains a factor q̇k q̇l . Retaining any terms in the
expansion for mkl (q) beyond the constant mkl (0) would lead to third and higher order terms.
We therefore keep only the constant term mkl (0). Thus, for small oscillations, we can write
the kinetic energy as
1X
K(q̇) = mkl q̇k q̇l , (B.50)
2 k,l
where in writing mkl , we mean mkl (0), here and hence forward. This equation serves as the
definition of mkl which, similar to the definition of kkl by Eq. (B.47), can be written as
∂2 K
!
mkl ≡ , (B.51)
∂q̇k ∂q̇l 0
where the kinetic energy used in this equation is the form containing terms only up to
quadratic order given by Eq. (B.50). Note that mkl = mlk as the order of differentiation
doesn’t matter.
The Lagrangian of a linearized system of coupled oscillators is thus given by
L(q, q̇) = K(q̇) − U(q) (B.52)
1X 1X
= mkl q̇k q̇l − kkl qk ql . (B.53)
2 k,l 2 k,l
We can obtain the generalized equations of motion for small oscillations of a system of
coupled oscillators from Lagrange’s equations. Because the kinetic energy is a function
only or the generalized velocities q̇ and the potential energy is a function only or the gen-
eralized coordinates q Lagrange’s equations become
d ∂K ∂U
!
= . (B.54)
dt ∂q̇k ∂qk
Using Eqs. (B.48) and (B.50) for K and U gives
X X
mkl q̈k = − kkl qk , (B.55)
k k
which can be written in matrix form

m q̈ = −k q . (B.56)
241 Problems
t
Here the coordinates q are understood to be generalized coordinates and not the normal
coordinates. Note that the matrices m and k are symmetric owing to their definitions by
Eqs. (B.51) and (B.47). You can use the fact that m and k are symmetric to check that
you have not made a mistake in setting up the equations of motion for a linearized system
of coupled oscillators, whether or not you use the Lagrangian formulation to obtain the
equations of motion.
Problems
B.1 Consider the problem of a pendulum swinging from an oscillating mass worked in
§??. Show that the Lagrangian of the system is given by
1 1 1
L= (M + m) ẋ2 + ml ẋ θ̇ cos θ + ml2 θ̇2 − mgl(1 − cos θ) − k ẋ2 . (B.57)
2 2 2
Use Eq. (B.22), Lagrange’s equations, to find the equations of motion of the system
and show that you get the same results as those obtained directly from Newton’s 2nd
law, which are given in §?? by Eqs. (??) and (??).
C Appendix C Computer programs
Here we go again ...
C.1 Jacobi method
...
242
Notes
243
245 Notes
t
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Oscillations & Waves

An intermediate introduction to the physics of waves

List of illustrations page viii

5 Waves of oscillating particles 129

6 Fourier Analysis 166

8 Sound and Electromagnetic Waves 217

9 Interference, Diffraction, and Fourier Optics 219

Appendix A Linear algebra 223

Appendix B Lagrangian formulation of classical mechanics 232

Appendix C Computer programs 242

1.1 Simple pendulum 4

1.1 The simple pendulum

the arc of the pendulum’s path:

θ(t) = a cos ω0 t + b sin ω0 t . (1.6)

θ(0) ≡ θ0 = a cos(ω0 · 0) + b sin(ω0 · 0) = a (1.7)

where we have introduced the notation θ̇ ≡ dθ/dt = Ω(t). Thus

θ(t) = A cos(ω0 t − δ) (1.10)

Fr ' −mg θ − 3!1 θ3 + 5!1 θ5 − ...

1.2 Diatomic molecule

1.3 The ideal spring

1.3.1 The horizontal spring

1.3.2 The vertical spring: an inhomogeneous ODE

x(t) = x1 + z(t) = l0 − (mg/k) sin θ (1 + cos ω0 t).

The maximum displacement is thus x(T/2) = l0 .

1.3.3 Initial conditions and constraints

where x = 0 is the initial and equilibrium position of mass M.

1.4 Energy of a simple harmonic oscillator

calculate forces by differentiation. In the second, we calculate force directly, remembering

(b) For motion along the y-direction, l1 = l2 = (2l0 )2 + y2 .

1.5 Complex exponential notation

Simple harmonic oscillations are described by sinusoidal functions. Unfortunately, work-

1.5.1 Complex exponentials

1.5.2 Complex exponentials in harmonic oscillator problems

sinh x = 12 e x − e−x (1.80)

Show that sinh(x) = −i sin(ix) and that cosh(x) = cos(ix).

y(t) = Re a1 eiωt + a2 e−iωt ,

y(t) = A1 cos ωt + A2 sin ωt , (1.83)

1.6 Series expansions

Example 1.3 Find the Taylor series expansion of f (y) = 1/ a2 − y2 .

1.7 Summary of important points of Chapter 1

for the system.

U(x) = V0 (1 − cos αx) − mgx.

2.1 The damped oscillator

where Γ is an experimentally-determined system-dependent damping or friction constant.

from the equation of motion, to obtain

−ω2 − iγω + ω20 = 0 . (2.13)

2.1.1 Underdamped response: (γ/2)2 < ω20

where ωu is a positive real frequency given by

zu1 (t) = a1 e−iω+ t = a1 e−γt/2 eiωu t (2.17)

1 −γt/2  i(ωu t−δ)

zu (t) = e−γt/2 (A1 cos ωu t + A2 sin ωu t) , (2.23)

oscillator can be written. We summarize them here:

The solutions correspond to different ways of expressing a sinusoidal function with an

The quality factor of an oscillator

2.1.2 Overdamped response: (γ/2)2 > ω20

into Eq. (2.32) gives1

2.1.3 Critically damped oscillations: (γ/2)2 = ω20

zc (t) = 30 t e−γt/2 . (2.44)

2.1.4 Summary of damped oscillations

The instantaneous rate of energy loss is obtained by differentiating Eq. (2.47)

3(t) = ż(t) ' 30 e−γt/2 cos ω0 t . (2.52)

2.3 Shock absorbers

ξc (t) = ξ̇i t e−γt/2 . (2.72)

ξ (t) = z(t) − z0 (m)

1 −γt/2 i(ωu t−δ)

Fig. 4.4 Displacements of two pendulums weakly coupled