INTRODUCTION

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 CHAPTER-1
INTRODUCTION
An automatic, switch activated paper-counting machine was designed to enable people
with disabilities to more easily perform task related to paper handling, such as counting papers
and separating papers. Our paper counting machine’s principle is based on the working of
currency counting machine and photo copying machine. The earlier machine is made for only
counting the pages but in large scale. The machine is counting pages and separating from that
bunch. In machine user enter number that much paper is separated from bunch. In that machine
we used DC motor for the rolling the papers and separate them out of stack. We are using keypad
for giving input of pages as per requirement. In that machine whenever we give input, it counts
pages. In that microcontroller interact with DC motor Driver, keypad, LCD, LDR sensor.

1.1 CURRENT SYSTEM:

The paper counting Machine & separating is advanced in machine basically its principle
is based on currency counting machine and photo copier machine. At present in most of the
places, counting is carried by manually so the probability of mistake is more. In the present
system most of the colleges and stationary stores, separating of pages is carried out manually.
After survey we found that only some of the industries who uses this type of machine for
counting and separating of pages in small scale. Julia O’Rourke, Peter Doblar, Jacob Felkl,
Mukund Kumar, Patrick Pace (University of Texas at Austin) these guys create a machine which
counts only pages. In advanced we proposed a new machine which counts as well as sort out
pages.
1.2 PROPOSED SYSTEM:

In previous machine only counting of pages carried out. In day to day life counting and
separating of pages is manually normally in our stationary shop counting is done by manually.
The currency counting machine & photo copier machine only pages are counted, to overcome
this drawback we design a small scale used machine which is small in dimension and reliable
than current paper separating system. In this machine both counting & separating of pages are
simultaneously carried out easily. Following demonstrate the Paper counting & separating
machine.
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1.3 METHOD OF COUNTING
1.3.1 CENTRE COUNTING:

In centre counting using two rotors are used to count the pages. Rotation of rotors is
given by using DC motor. Rotor is placed at the centre position of the paper. In this method the
DC motor are rotating in clockwise manner it sort out the pages these pages are collected in tray.
In this method of centre counting mechanism one drawback is it can't count pages normal
method.

1.3.2 LASER COUNTING:

The Laser counting mechanism is normally used in currency counting machine or
printing machine. In these machine laser gun is interfaced with microcontroller when pages cut
the light beam and received by receiver and gives output to the microcontroller. In these
mechanisms we count pages. For separating process we used a steeper motor to remove pages.

1.3.3 CORNER COUNTING:

In corner counting rotors are used to count the pages. In this mechanism the rotors are
placed in top corner side of machine using this method we count as well as sorts the pages. When
rotors are rotating in clockwise direction then pages are one by one collect in bottom side tray.
Rotors rotate with the help of DC motor the pages counting process is start.

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COMPONENT
DESCRIPTION

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 CHAPTER-2
COMPONENT DESCRIPTION

The paper counting machine consist the following components as it important parts,

Those are,

 MICROCONTROLLER (AT89C51)
 LCD DISPLAY
 DC MOTOR
 SENSOR
 RELAYS
 TRANSFORMER
 DIODES
 RESISTORS
 CAPACITORS
 INPUT TRAY
 OUTPUT TRAY
 ROLLER

CHAPTER-3
BRIEF EXPLANATION OF COMPONENTS

3.1 AT89C51 MICROCONTROLLER

The AT89C51 is a low-power, high-performance CMOS 8-bit microcomputer with 4K

bytes of Flash Programmable and Erasable Read Only Memory (PEROM). The device is

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manufactured using Atmel’s high density non-volatile memory technology and is compatible
with the industry standard MCS-51Ô instruction set and pin out. The on-chip Flash allows the
program memory to be reprogrammed in-system or by a conventional non-volatile memory
programmer. By combining a versatile 8-bit CPU with Flash on a monolithic chip, the Atmel
AT89C51 is a powerful microcomputer which provides a highly flexible and cost effective
solution to many embedded control applications.

The Intel MCS-51 (commonly referred to as 8051) is Harvard architecture, single chip
microcontroller (µC) series which was developed by Intel in 1980 for use in embedded systems.
Intel's original versions were popular in the 1980s and early 1990s. While Intel no longer
manufactures the MCS-51, binary compatible derivatives remain popular today. In addition to
these physical devices, several companies also offer MCS-51 derivatives as IP cores for use in
FPGAs or ASICs designs.
Intel's original MCS-51 family was developed using NMOS technology, but later
versions, identified by a letter C in their name (e.g., 80C51) used CMOS technology and
consumed less power than their NMOS predecessors. This made them more suitable for battery-
powered devices.

AT89C51 Microcontroller.
3.2 AT89C51

Important features and applications of 8051 micro architecture.

 The 8051 architecture provides many functions (CPU, RAM, ROM, I/O, interrupt logic,
timer, etc.) in a single package
 8-bit ALU, Accumulator and 8-bit Registers; hence it is an 8-bit microcontroller
 8-bit data bus – It can access 8 bits of data in one operation

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  16-bit address bus – It can access 216 memory locations – 64 KB (65536 locations) each
of RAM and ROM
 On-chip RAM – 128 bytes (data memory)
 On-chip ROM – 4 Kbyte (program memory)
 Four byte bi-directional input/output port
 UART (serial port)
 Two 16-bit Counter/timers
 Two-level interrupt priority
 Power saving mode (on some derivatives)
One particularly useful feature of the 8051 core was the inclusion of a Boolean
processing engine which allows bit-level Boolean logic operations to be carried out directly and
efficiently on select internal registers and select RAM locations. This advantageous feature
helped cement the 8051's popularity in industrial control applications because it reduced code
size by as much as 30%. Another valued feature is the including of four bank selectable working
register sets which greatly reduce the amount of time required to complete an interrupt service
routine.

With a single instruction the 8051 can switch register banks as opposed to the time
consuming task of transferring the critical registers to the stack or designated RAM locations.
These registers also allowed the 8051 to quickly perform a context switch which is essential for
time sensitive real-time applications. The MCS-51 UARTs make it simple to use the chip as a
serial communications interface. External pins can be configured to connect to internal shift
registers in a variety of ways, and the internal timers can also be used, allowing serial
communications in a number of modes, both synchronous and asynchronous. Some modes allow
communications with no external components.

A mode compatible with an RS-485 multi-point communications environment is

achievable, but the 8051's real strength is fitting in with existing ad-hoc protocols (e.g., when
controlling serial-controlled devices).Once a UART, and a timer if necessary, have been
configured, the programmer needs only to write a simple interrupt routine to refill the send shift
register whenever the last bit is shifted out by the UART and/or empty the full receive shift
register (copy the data somewhere else). The main program then performs serial reads and writes
simply by reading and writing 8-bit data to stacks.

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 MCS-51 based microcontrollers typically include one or two UARTs, two or three timers,
128 or 256 bytes of internal data RAM (16 bytes of which are bit-addressable), up to 128 bytes
of I/O, 512 bytes to 64 kB of internal program memory, and sometimes a quantity of extended
data RAM (ERAM) located in the external data space. The original 8051 core ran at 12 clock
cycles per machine cycle, with most instructions executing in one or two machine cycles. With a
12 MHz clock frequency, the 8051 could thus execute 1 million one-cycle instructions per
second or 500,000 two-cycle instructions per second.

Enhanced 8051 cores are now commonly used which run at six, four, two, or even one
clock per machine cycle, and have clock frequencies of up to 100 MHz, and are thus capable of
an even greater number of instructions per second. All SI Labs, some Dallas and a few Atmel
devices have single cycle cores. Features of the modern 8051 include built-in reset timers with
brown-out detection, on-chip oscillators, self-programmable Flash ROM program memory, built-
in external RAM, extra internal program storage, boot loader code in ROM, EEPROM non-
volatile data storage, I²C, SPI, and USB host interfaces, CAN or LIN bus, PWM generators,
analog comparators, A/D and D/A converters, RTCs, extra counters and timers, in-circuit
debugging facilities, more interrupt sources, and extra power saving modes. In many engineering
schools the 8051 microcontroller is used in introductory microcontroller courses.

3.3 MEMORY ARCHITECTURE

The MCS-51 has four distinct types of memory – internal RAM, special function
registers, program memory, and external data memory. Internal RAM (IRAM) is located from
address 0 to address 0xFF. IRAM from 0x00 to 0x7F can be accessed directly, and the bytes from
0x20 to 0x2F are also bit-addressable. IRAM from 0x80 to 0xFF must be accessed indirectly,
using the @R0 or @R1 syntax, with the address to access loaded in R0 or R1.

Special function registers (SFR) are located from address 0x80 to 0xFF, and are accessed
directly using the same instructions as for the lower half of IRAM. Some of the SFR's are also
bit-addressable. Program memory (PMEM, though less common in usage than IRAM and

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XRAM) is located starting at address 0. It may be on- or off-chip, depending on the particular
model of chip being used. Program memory is read-only, though some variants of the 8051 use
on-chip flash memory and provide a method of re-programming the memory in-system or in-
application. Aside from storing code, program memory can also store tables of constants that can
be accessed by MOVC A, @DPTR, using the 16-bit special function register DPTR.

External data memory (XRAM) also starts at address 0. It can also be on- or off-chip;
what makes it "external" is that it must be accessed using the MOVX (Move external)
instruction. Many variants of the 8051 include the standard 256 bytes of IRAM plus a few KB of
XRAM on the chip. If more XRAM is required by an application, the internal XRAM can be
disabled, and all MOVX instructions will fetch from the external bus.

INTEL 8051 block diagram.

3.4 RELAY

A relay is an electrically operated switch. Many relays use an electromagnet to operate a

switching mechanism mechanically, but other operating principles are also used. Relays are used
where it is necessary to control a circuit by a low-power signal (with complete electrical isolation
between control and controlled circuits), or where several circuits must be controlled by one
signal. The first relays were used in long distance telegraph circuits, repeating the signal coming

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in from one circuit and re-transmitting it to another. Relays were used extensively in telephone
exchanges and early computers to perform logical operations.

A type of relay that can handle the high power required to directly control an electric
motor or other loads is called a contractor. Solid-state relays control power circuits with no
moving parts, instead using a semiconductor device to perform switching. Relays with calibrated
operating characteristics and sometimes multiple operating coils are used to protect electrical
circuits from overload or faults; in modern electric power systems these functions are performed
by digital instruments still called "protective relays".

Relay is a common, simple application of electromagnetism. It uses an electromagnet

made from an iron rod wound with hundreds of fine copper wire. When electricity is applied to
the wire, the rod becomes magnetic. A movable contact arm above the rod is then pulled toward
the rod until it closes a switch contact. When the electricity is removed, a small spring pulls the
contract arm away from the rod until it closes a second switch contact. By means of relay, a
current circuit can be broken or closed in one circuit as a result of a current in another circuit.

3.5 BASIC DESIGN AND OPERATION

Small "cradle" relay often used in electronics. The "cradle" term refers to the shape of the
relay's armature. A simple electromagnetic relay consists of a coil of wire wrapped around a soft
iron core, an iron yoke which provides a low reluctance path for magnetic flux, a movable iron
armature, and one or more sets of contacts (there are two in the relay pictured). The armature is
hinged to the yoke and mechanically linked to one or more sets of moving contacts. It is held in
place by a spring so that when the relay is de-energized there is an air gap in the magnetic circuit.
In this condition, one of the two sets of contacts in the relay pictured is closed, and the other set
is open. Other relays may have more or fewer sets of contacts depending on their function. The
relay in the picture also has a wire connecting the armature to the yoke. This ensures continuity
of the circuit between the moving contacts on the armature, and the circuit track on the printed
circuit board (PCB) via the yoke, which is soldered to the PCB.

When an electric current is passed through the coil it generates a magnetic field that
activates the armature and the consequent movement of the movable contact either makes or
breaks (depending upon construction) a connection with a fixed contact. If the set of contacts

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was closed when the relay was de-energized, then the movement opens the contacts and breaks
the connection, and vice versa if the contacts were open. When the current to the coil is switched
off, the armature is returned by a force, approximately half as strong as the magnetic force, to its
relaxed position. Usually this force is provided by a spring, but gravity is also used commonly in
industrial motor starters. Most relays are manufactured to operate quickly. In a low-voltage
application this reduces noise; in a high voltage or current application it reduces arcing.

When the coil is energized with direct current, a diode is often placed across the coil to
dissipate the energy from the collapsing magnetic field at deactivation, which would otherwise
generate a voltage spike dangerous to semiconductor circuit components. Some automotive
relays include a diode inside the relay case.

Alternatively, a contact protection network consisting of a capacitor and resistor in series

(snubber circuit) may absorb the surge. If the coil is designed to be energized with alternating
current (AC), a small copper "shading ring" can be crimped to the end of the solenoid, creating a
small out-of-phase current which increases the minimum pull on the armature during the AC
cycle.A solid-state relay uses a thyristor or other solid-state switching device, activated by the
control signal, to switch the controlled load, instead of a solenoid. An optocoupler (a light-
emitting diode (LED) coupled with a photo transistor) can be used to isolate control and
controlled circuits.
3.6 RELAY SWITCH

FIG- 7 Relay switch

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3.7 CRYSTAL OSCILLATORS
A crystal oscillator is an electronic oscillator circuit that uses the mechanical resonance of
a vibrating crystal of piezoelectric material to create an electrical signal with a very precise
frequency. This frequency is commonly used to keep track of time (as in quartz wristwatches), to
provide a stable clock signal for digital integrated circuits, and to stabilize frequencies for radio
transmitters and receivers. The most common type of piezoelectric resonator used is the quartz
crystal, so oscillator circuits designed around them became known as "crystal oscillators."
Quartz crystals are manufactured for frequencies from a few tens of kilohertz to tens of
megahertz. More than two billion (2×109) crystals are manufactured annually. Most are used for
consumer devices such as wristwatches, clocks, radios, computers, and cell phones. Quartz
crystals are also found inside test and measurement equipment, such as counters, signal
generators, and oscilloscopes.

Crystal oscillators are oscillators where the primary frequency determining element is a
quartz crystal. Because of the inherent characteristics of the quartz crystal the crystal oscillator
may be held to extreme accuracy of frequency stability. Temperature compensation may be
applied to crystal oscillators to improve thermal stability of the crystal oscillator. Crystal
oscillators are usually, fixed frequency oscillators where stability and accuracy are the primary
considerations. For example it is almost impossible to design a stable and accurate LC oscillator
for the upper HF and higher frequencies without resorting to some sort of crystal control.

3.8 OPERATION

A crystal is a solid in which the constituent atoms, molecules, or ions are packed in a
regularly ordered, repeating pattern extending in all three spatial dimensions.
Almost any object made of an elastic material could be used like a crystal, with appropriate
transducers, since all objects have natural resonant frequencies of vibration. For example, steel is
very elastic and has a high speed of sound. It was often used in mechanical filters before quartz.
The resonant frequency depends on size, shape, elasticity, and the speed of sound in the material.
High-frequency crystals are typically cut in the shape of a simple, rectangular plate. Low-
frequency crystals, such as those used in digital watches, are typically cut in the shape of a
tuning fork.

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 For applications not needing very precise timing, a low-cost ceramic resonator is often
used in place of a quartz crystal. When a crystal of quartz is properly cut and mounted, it can be
made to distort in an electric field by applying a voltage to an electrode near or on the crystal.
This property is known as piezoelectricity. When the field is removed, the quartz will generate an
electric field as it returns to its previous shape, and this can generate a voltage. The result is that
a quartz crystal behaves like a circuit composed of an inductor, capacitor and resistor, with a
precise resonant frequency.
Quartz has the further advantage that its elastic constants and its size change in such a
way that the frequency dependence on temperature can be very low. The specific characteristics
will depend on the mode of vibration and the angle at which the quartz is cut (relative to its
crystallographic axes). Therefore, the resonant frequency of the plate, which depends on its size,
will not change much, either.
This means that a quartz clock, filter or oscillator will remain accurate. For critical applications
the quartz oscillator is mounted in a temperature-controlled container, called a crystal oven, and
can also be mounted on shock absorbers to prevent perturbation by external mechanical
vibrations.

3.9 CRYSTAL STRUCTURES AND MATERIALS

The most common material for oscillator crystals is quartz. At the beginning of the
technology, natural quartz crystals were used; now synthetic crystalline quartz grown by
hydrothermal synthesis is predominant due to higher purity, lower cost, and more convenient
handling. One of the few remaining uses of natural crystals is for pressure transducers in deep
wells. During World War II and for some time afterwards, natural quartz was considered a
strategic material by the USA. Large crystals were imported from Brazil. Raw "lascas", the
source material quartz for hydrothermal synthesis, are imported to USA or mined locally by
Coleman Quartz. The average value of as-grown synthetic quartz in 1994 was USD60/kg.

Two types of quartz crystals exist: left-handed and right-handed, differing in the optical
rotation but identical in other physical properties. Both left and right-handed crystals can be used
for oscillators, if the cut angle is correct. In manufacture, right-handed quartz is generally used.

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The SiO4 tetrahedrons form parallel helixes; the direction of twist of the helix determines the
left- or right-hand orientation. The helixes are aligned along the z-axis and merged together,
sharing atoms. The mass of the helixes forms a mesh of small and large channels parallel to the
z-axis; the large ones are large enough to allow some mobility of smaller ions and molecules
through the crystal.

Quartz exists in several phases. At 573 °C at 1 atmosphere (and at higher temperatures

and higher pressures) the α-quartz undergoes quartz inversion, transforms reversibly to β-quartz.
The reverse process however is not entirely homogeneous and crystal twinning occurs. Care has
to be taken during manufacture and processing to avoid the phase transformation. Other phases,
e.g. the higher-temperature phases tridymite and cristobalite, are not significant for oscillators.
All quartz oscillator crystals are the α-quartz type.

Infrared spectro photometry is used as one of the methods for measuring the quality of
the grown crystals. The wave numbers 3585, 3500 and 3410 cm−1 are commonly used. The
measured value is based on the absorption bands of the OH radical and the infrared Q value is
calculated. The electronic grade crystals, grade C, have Q of 1.8 million or above; the premium
grade B crystals have Q of 2.2 million, and special premium grade A crystals have Q of 3.0
million. The Q value is calculated only for the z region; crystals containing other regions can be
adversely affected. Another quality indicator is the etch channel density; when the crystal is
etched, tubular channels are created along linear defects. For processing involving etching, e.g.
the wristwatch tuning fork crystals, low etch channel density is desirable. The etch channel
density for swept quartz is about 10–100 and significantly more for unwept quartz. Presence of
etch channels and etch pits degrades the resonator's Q and introduces nonlinearities. Quartz
crystals can be grown for specific purposes.

Crystals for AT-cut are the most common in mass production of oscillator materials; the
shape and dimensions are optimized for high yield of the required wafers. High-purity quartz
crystals are grown with especially low content of aluminium, alkali metal and other impurities
and minimal defects; the low amount of alkali metals provides increased resistance to ionizing
radiation. Crystals for wrist watches, for cutting the tuning fork 32768 Hz crystals, are grown

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with very low etch channel density. Crystals for SAW devices are grown as flat; with large X-
size seed with low etch channel density.

Special high-Q crystals, for use in highly stable oscillators, are grown at constant slow
speed and have constant low infrared absorption along the entire Z axis. Crystals can be grown
as Y-bar, with a seed crystal in bar shape and elongated along the Y axis, or as Z-plate, grown
from a plate seed with Y-axis direction length and X-axis width. The region around the seed
crystal contains a large number of crystal defects and should not be used for the wafers. Crystals
grow anisotropically; the growth along the Z axis is up to 3 times faster than along the X axis.
The growth direction and rate also influences the rate of uptake of impurities. Y-bar crystals or Z-
plate crystals with long Y axis, have four growth regions usually called + X, -X, Z, and S. The
distribution of impurities during growth is uneven; different growth areas contain different level
of contaminants. The z regions are the purest, the small occasionally present s regions are less
pure, the +x region is yet less pure, and the -x region has the highest level of impurities.

The impurities have negative impact on radiation hardness, susceptibility to twinning,

filter loss, and long and short term stability of the crystals. Different-cut seeds in different
orientations may provide other kinds of growth regions. The growth speed of the -x direction is
slowest due to the effect of adsorption of water molecules on the crystal surface; aluminium
impurities suppress growth in two other directions. The content of aluminium is lowest in z
region, higher in +x, yet higher in -x, and highest in s; the size of s regions also grows with
increased amount of aluminium present.

The content of hydrogen is lowest in z region, higher in +x region, yet higher in s region,
and highest in -x. Aluminium inclusions transform to colour centres with a gamma ray
irradiation, causing darkening of the crystal proportional to the dose and level of impurities;
presence of regions with different darkness reveals the different growth regions. The dominant
type of defect of concern in quartz crystals is the substitution of Al(III) for Si(IV) atom in the
crystal lattice. The aluminium ion has an associated interstitial charge compensator present
nearby, which can be a H+ ion (attached to the nearby oxygen and forming a hydroxyl group,
called Al-OH defect), Li+ ion, Na+ ion, K+ ion (less common), or an electron hole trapped in a
nearby oxygen atom orbital. The composition of the growth solution, whether it is based on

15
lithium or sodium alkali compounds, determines the charge compensating ions for the aluminium
defects. The ion impurities are of concern as they are not firmly bound and can migrate through
the crystal, altering the local lattice elasticity and the resonant frequency of the crystal. Other
common impurities of concern are e.g. iron(III) (interstitial), fluorine, boron(III), phosphorus(V)
(substitution), titanium(IV) (substitution, universally present in magmatic quartz, less common in
hydrothermal quartz), and germanium(IV) (substitution).

Sodium and iron ions can cause inclusions of aconite and elemeusite crystals. Inclusions
of water may be present in fast-grown crystals; interstitial water molecules are abundant near the
crystal seed. Another defect of importance is the hydrogen containing growth defect, when
instead of a Si-O-Si structure a pair of Si-OH HO-Si groups is formed; essentially a hydrolyzed
bond. Fast-grown crystals contain more hydrogen defects than slow-grown ones. These growth
defects source as supply of hydrogen ions for radiation-induced processes and forming Al-OH
defects. Germanium impurities tend to trap electrons created during irradiation; the alkali metal
cations then migrate towards the negatively charged centre and form a stabilizing complex.
Matrix defects can be also present; oxygen vacancies, silicon vacancies (usually compensated by
4 hydrogen or 3 hydrogen and a hole), peroxy groups, etc. Some of the defects produce localized
levels in the forbidden band, serving as charge traps; Al (III) and B (III) typically serve as hole
traps while electron vacancies, titanium, germanium, and phosphorus atoms serve as electron
traps. The trapped charge carriers can be released by heating; their recombination is the cause of
thermo luminescence.

The mobility of interstitial ions depends strongly on temperature. Hydrogen ions are
mobile down to 10 K, but alkali metal ions become mobile only at temperatures around and
above 200 K. The hydroxyl defects can be measured by near-infrared spectroscopy. The trapped
holes can be measured by electron spin resonance. The Al-Na+ defects show as an acoustic loss
peak due to their stress-induced motion; the Al-Li+ defects do not form a potential well so are
not detectable this way. Some of the radiation induced defects during their thermal annealing
produce thermo luminescence; defects related to aluminium, titanium, and germanium can be
distinguished.

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 Swept crystals are crystals that have undergone a solid-state electro diffusion purification
process. Sweeping involves heating the crystal above 500 °C in a hydrogen-free atmosphere, and
the voltage gradient of at least 1 kilovolt/cm, for several (usually over 12) hours. The migration
of impurities and the gradual replacement of alkali metal ions with hydrogen (when swept in air)
or electron holes (when swept in vacuum) causes a weak electric current through the crystal;
decay of this current to a constant value signals end of the process. The crystal is then left to
cool, while the electric field is maintained. The impurities are concentrated at the cathode region
of the crystal, which is cut off afterwards and discarded. Swept crystals have increased resistance
to radiation, as the dose effects are dependent on the level of alkali metal impurities; they are
suitable for use in devices exposed to ionizing radiation, e.g. for nuclear and space technology.
Sweeping under vacuum at higher temperatures and higher field strengths yields yet more
radiation-hard crystals. The level and character of impurities can be measured by infrared
spectroscopy. Quartz can be swept in both α and β phase; sweeping in β phase is faster, but the
phase transition may induce twinning. Twinning can be mitigated by subjecting the crystal to
compression stress in the X direction, or an AC or DC electric field along the X axis while the
crystal cools through the phase transformation temperature region.

Sweeping can be also used to introduce one kind of an impurity into the crystal. Lithium,
sodium, and hydrogen swept crystals are used for e.g. studying quartz behavior.Very small
crystals for high fundamental mode frequencies can be manufactured by photolithography.
Crystals can be adjusted to exact frequency by laser trimming. A technique used in the world of
amateur radio for slight decrease of the crystal frequency may be achieved by exposing crystals
with silver electrodes to vapours of iodine, which causes a slight mass increase on the surface by
forming a thin layer of silver iodide; such crystals however had problematic long-term stability.
Another method commonly used is electrochemical increase or decrease of silver electrode
thickness by submerging resonator in lapis solved in water, citric acid in water, or water with
salt, and using resonator as one electrode, and small silver electrode as another.

By choosing direction of current, one can either increase or decrease mass of electrodes. Details
were published in "Radio" magazine (3/1978) by UB5LEV.Raising frequency by scratching off
parts of the electrodes is advised against, as this may damage the crystal and lower its Q factor.
Capacitor trimmers can be also used for frequency adjustment of the oscillator circuit.

17
 Some other piezoelectric materials than quartz can be employed; e.g. single crystals of
lithium tantalite, lithium niobate, lithium borate, berlinite, gallium arsenide, lithium tetraborate,
aluminium phosphate, bismuth germanium oxide, polycrystalline zirconium titanate ceramics,
high-alumina ceramics, silicon-zinc oxide composite, or dipotassium tartrate; some materials
may be more suitable for specific applications. An oscillator crystal can be also manufactured by
depositing the resonator material on the silicon chip surface. Crystals of gallium phosphate,
langasite, langanite and langanate are about 10 times more pullable than the corresponding
quartz crystals, and are used in some VCXO oscillators.

3.10 RESISTORS

Resistors (R), are the most commonly used of all electronic components, to the point
where they are almost taken for granted. There are many different resistor types available with
their principal job being to "resist" the flow of current through an electrical circuit, or to act as
voltage droppers or voltage dividers. They are "Passive Devices", that is they contain no source
of power or amplification but only attenuate or reduce the voltage signal passing through them.
When used in DC circuits the voltage drop produced is measured across their terminals as the
circuit current flows through them while in AC circuits the voltage and current are both in-
phase producing 0o phase shift.

Resistors produce a voltage drop across themselves when an electrical current flows
through them because they obey Ohm's Law, and different values of resistance produces different
values of current or voltage. This can be very useful in Electronic circuits by controlling or
reducing either the current flow or voltage produced across them. There are many
different Resistor Types and they are produced in a variety of forms because their particular
characteristics and accuracy suit certain areas of application, such as High Stability, High
Voltage, High Current etc., or are used as general purpose resistors where their characteristics are
less of a problem. Some of the common characteristics associated with the humble resistor

18
are; Temperature Coefficient, Voltage Coefficient, Noise, Frequency Response, Power as well
as Temperature Rating, Physical Size and Reliability.

In all Electrical and Electronic circuit diagrams and schematics, the most commonly used
resistor symbol is that of a "zigzag" type line with the value of its resistance given in Ohms, Ω.

3.11 CAPACITOR

Just like the Resistor, the Capacitor or sometimes referred to as a Condenser is a passive
device, and one which stores energy in the form of an electrostatic field which produces a
potential (Static Voltage) across its plates. In its basic form a capacitor consists of two parallel
conductive plates that are not connected but are electrically separated either by air or by an
insulating material called the Dielectric. When a voltage is applied to these plates, a current
flows charging up the plates with electrons giving one plate a positive charge and the other plate
an equal and opposite negative charge. This flow of electrons to the plates is known as
the Charging Current and continues to flow until the voltage across the plates (and hence the
capacitor) is equal to the applied voltage Vc. At this point the capacitor is said to be fully charged
and this is illustrated below.

3.12 CAPACITOR CONSTRUCTION

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 FIG- 8 :Capacitor construction

The parallel plate capacitor is the simplest form of capacitor and its capacitance value is
fixed by the equal area of the plates and the distance or separation between them. Altering any
two of these values alters the value of its capacitance and this forms the basis of operation of the
variable capacitors. Also, because capacitors store the energy of the electrons in the form of an
electrical charge on the plates the larger the plates and/or smaller their separation the greater will
be the charge that the capacitor holds for any given voltage across its plates.

3.13 LIQUID CRYSTAL DISPLAY

A liquid crystal display (LCD) is a flat panel display, electronic visual display, or video
display that uses the light modulating properties of liquid crystals (LCs). LCs do not emit light
directly.

20
 A general purpose alphanumeric LCD, with two lines of 16 characters.

LCDs are used in a wide range of applications, including computer monitors, television,
instrument panels, aircraft cockpit displays, signage, etc. They are common in consumer devices
such as video players, gaming devices, clocks, watches, calculators, and telephones. LCDs have
replaced cathode ray tube (CRT) displays in most applications. They are available in a wider
range of screen sizes than CRT and plasma displays, and since they do not use phosphors, they
cannot suffer image burn-in. LCDs are, however, susceptible to image persistence.
LCDs are more energy efficient and offer safer disposal than CRTs. Its low electrical
power consumption enables it to be used in battery-powered electronic equipment. It is an
electronically modulated optical device made up of any number of segments filled with liquid
crystals and arrayed in front of a light source (backlight) or reflector to produce images in color
or monochrome. The most flexible ones use an array of small pixels. The earliest discovery
leading to the development of LCD technology, the discovery of liquid crystals, dates from 1888.
By 2008, worldwide sales of televisions with LCD screens had surpassed the sale of CRT units.

3.14 OVERVIEW

Each pixel of an LCD typically consists of a layer of molecules aligned between two
transparent electrodes, and two polarizing filters, the axes of transmission of which are (in most
of the cases) perpendicular to each other. With no actual liquid crystal between the polarizing
filters, light passing through the first filter would be blocked by the second (crossed) polarizer.

21
 The surfaces of the electrodes that are in contact with the liquid crystal material are
treated so as to align the liquid crystal molecules in a particular direction. This treatment
typically consists of a thin polymer layer that is unidirectional rubbed using, for example, a cloth.
The direction of the liquid crystal alignment is then defined by the direction of rubbing.
Electrodes are made of the transparent conductor Indium Tin Oxide (ITO). The Liquid Crystal
Display is intrinsically a “passive” device; it is a simple light valve. The managing and control of
the data to be displayed is performed by one or more circuits commonly denoted as LCD drivers.

Before applying an electric field, the orientation of the liquid crystal molecules is
determined by the alignment at the surfaces of electrodes. In a twisted nematic device (still the
most common liquid crystal device), the surface alignment directions at the two electrodes are
perpendicular to each other, and so the molecules arrange themselves in a helical structure, or
twist. This induces the rotation of the polarization of the incident light, and the device appears
grey. If the applied voltage is large enough, the liquid crystal molecules in the center of the layer
are almost completely untwisted and the polarization of the incident light is not rotated as it
passes through the liquid crystal layer. This light will then be mainly polarized perpendicular to
the second filter, and thus be blocked and the pixel will appear black. By controlling the voltage
applied across the liquid crystal layer in each pixel, light can be allowed to pass through in
varying amounts thus constituting different levels of gray.

3.15 Voltage Regulator

The 78xx (sometimes LM78xx) is a family of self-contained fixed linear voltage

regulator integrated circuits. The 78xx family is commonly used in electronic circuits requiring a
regulated power supply due to their ease-of-use and low cost. For ICs within the family, the xx is
replaced with two digits, indicating the output voltage (for example, the 7805 has a 5 volt output,
while the 7812 produces 12 volts). The 78xx lines are positive voltage regulators: they produce a
voltage that is positive relative to a common ground. There is a related line of 79xx devices
which are complementary negative voltage regulators. 78xx and 79xx ICs can be used in
combination to provide positive and negative supply voltages in the same circuit.

When the input voltage is significantly higher than the regulated output voltage (for
example, powering a 7805 using a 24 volt power source), this inefficiency can be a significant

22
issue. Even in larger packages, 78xx integrated circuits cannot supply as much power as many
designs which use discrete components, and are generally inappropriate for applications
requiring more than a few amperes of current.

3.16 TRANSFORMER

A transformer is a device that transfers electrical energy from one circuit to another
through inductively coupled conductors—the transformer's coils. A varying current in the first or
primary winding creates a varying magnetic flux in the transformer's core and thus a varying
magnetic field through the secondary winding. This varying magnetic field induces a varying
electromotive force (EMF), or "voltage", in the secondary winding. This effect is called inductive
coupling.

TRANSFORMER WINDINGS

If a load is connected to the secondary, current will flow in the secondary winding, and
electrical energy will be transferred from the primary circuit through the transformer to the load.
In an ideal transformer, the induced voltage in the secondary winding (Vs) is in proportion to the
primary voltage (Vp) and is given by the ratio of the number of turns in the secondary (Ns) to the
number of turns in the primary (Np) as follows: By appropriate selection of the ratio of turns, a
transformer thus enables an alternating current (AC) voltage to be "stepped up" by making Ns

23
greater than Np, or "stepped down" by making Ns less than Np.In the vast majority of
transformers, the windings are coils wound around a ferromagnetic core, air-core transformers
being a notable exception.

Transformers range in size from a thumbnail-sized coupling transformer hidden inside a

stage microphone to huge units weighing hundreds of tons used to interconnect portions of
power grids. All operate on the same basic principles, although the range of designs is wide.
While new technologies have eliminated the need for transformers in some electronic circuits,
transformers are still found in nearly all electronic devices designed for household ("mains")
voltage. Transformers are essential for high-voltage electric power transmission, which makes
long-distance transmission economically practical.

A transformer is an electrical device that transfers energy from one circuit to another by
magnetic coupling with no moving parts. A transformer comprises two or more coupled
windings, or a single tapped winding and, in most cases, a magnetic core to concentrate magnetic
flux. A changing current in one winding creates a time-varying magnetic flux in the core, which
induces a voltage in the other windings. Michael Faraday built the first transformer, although he
used it only to demonstrate the principle of electromagnetic induction and did not foresee the use
to which it would eventually be put.

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 Transformer core
3.17 DC MOTOR

A DC motor is an electric motor that runs on direct current (DC) electricity. DC motors
were used to run machinery, often eliminating the need for a local steam engine or internal
combustion engine. DC motors can operate directly from rechargeable batteries, providing the
motive power for the first electric vehicles. Today DC motors are still found in applications as
small as toys and disk drives, or in large sizes to operate steel rolling mills and paper machines.
Modern DC motors are nearly always operated in conjunction with power electronic devices.

Two important performance parameters of DC motors are the motor constants, Kv and
Km.

25
 DC motor

When a current passes through the coil wound around a soft iron core, the side of the
positive pole is acted upon by an upwards force, while the other side is acted upon by a
downward force. According to Fleming's left hand rule, the forces cause a turning effect on the
coil, making it rotate. To make the motor rotate in a constant direction, "direct current"
commutators make the current reverse in direction every half a cycle (in a two-pole motor) thus
causing the motor to continue to rotate in the same direction.

A problem with the motor shown above is that when the plane of the coil is parallel to the
magnetic field—i.e. when the rotor poles are 90 degrees from the stator poles—the torque is
zero. In the pictures above, this occurs when the core of the coil is horizontal—the position it is
just about to reach in the last picture on the right. The motor would not be able to start in this
position. However, once it was started, it would continue to rotate through this position by
momentum.

There is a second problem with this simple pole design. At the zero-torque position, both
commutator brushes are touching (bridging) both commutator plates, resulting in a short-circuit.
The power leads are shorted together through the commutator plates, and the coil is also short-
circuited through both brushes (the coil is shorted twice, once through each brush
independently). Note that this problem is independent of the non-starting problem above; even if

26
there were a high current in the coil at this position, there would still be zero torque. The problem
here is that this short uselessly consumes power without producing any motion (nor even any
coil current.) In a low-current battery-powered demonstration this short-circuiting is generally
not considered harmful. However, if a two-pole motor were designed to do actual work with
several hundred watts of power output, this shorting could result in severe commutator
overheating, brush damage, and potential welding of the brushes—if they were metallic—to the
commutator. Carbon brushes, which are often used, would not weld. In any case, a short like this
is very wasteful, drains batteries rapidly and, at a minimum, requires power supply components
to be designed to much higher standards than would be needed just to run the motor without the
shorting.

3.18 THE INSIDE OF AN ELECTRIC DC MOTOR.

One simple solution is to put a gap between the commutator plates which is wider than
the ends of the brushes. This increases the zero-torque range of angular positions but eliminates
the shorting problem; if the motor is started spinning by an outside force it will continue
spinning. With this modification, it can also be effectively turned off simply by stalling
(stopping) it in a position in the zero-torque (i.e. commutator non-contacting) angle range. This
design is sometimes seen in homebuilt hobby motors, e.g. for science fairs and such designs can
be found in some published science project books. A clear downside of this simple solution is
that the motor now coasts through a substantial arc of rotation twice per revolution and the
torque is pulsed. This may work for electric fans or to keep a flywheel spinning but there are
many applications, even where starting and stopping are not necessary, for which it is completely
inadequate, such as driving the capstan of a tape transport, or any instance where to speed up and
slow down often and quickly is a requirement. Another disadvantage is that, since the coils have
a measure of self inductance, current flowing in them cannot suddenly stop. The current attempts
to jump the opening gap between the commutator segment and the brush, causing arcing.

Even for fans and flywheels, the clear weaknesses remaining in this design—especially
that it is not self-starting from all positions—make it impractical for working use, especially
considering the better alternatives that exist. Unlike the demonstration motor above, DC motors
are commonly designed with more than two poles, are able to start from any position, and do not
have any position where current can flow without producing electromotive power by passing

27
through some coil. Many common small brushed DC motors used in toys and small consumer
appliances, the simplest mass-produced DC motors to be found, have three-pole armatures. The
brushes can now bridge two adjacent commutator segments without causing a short circuit.
These three-pole armatures also have the advantage that current from the brushes either flows
through two coils in series or through just one coil. Starting with the current in an individual coil
at half its nominal value (as a result of flowing through two coils in series), it rises to its nominal
value and then falls to half this value. The sequence then continues with current in the reverse
direction. This results in a closer step-wise approximation to the ideal sinusoidal coil current,
producing a more even torque than the two-pole motor where the current in each coil is closer to
a square wave. Since current changes are half those of a comparable two-pole motor, arcing at
the brushes is consequently less.

If the shaft of a DC motor is turned by an external force, the motor will act like a
generator and produce an Electromotive force (EMF). During normal operation, the spinning of
the motor produces a voltage, known as the counter-EMF (CEMF) or back EMF, because it
opposes the applied voltage on the motor. The back EMF is the reason that the motor when free-
running does not appear to have the same low electrical resistance as the wire contained in its
winding. This is the same EMF that is produced when the motor is used as a generator (for
example when an electrical load, such as a light bulb, is placed across the terminals of the motor
and the motor shaft is driven with an external torque). Therefore, the total voltage drop across a
motor consists of the CEMF voltage drop, and the parasitic voltage drop resulting from the
internal resistance of the armature's windings.

3.19 SPEED CONTROL

Generally, the rotational speed of a DC motor is proportional to the voltage applied to it,
and the torque is proportional to the current. Speed control can be achieved by variable battery
tappings, variable supply voltage, resistors or electronic controls. The direction of a wound field
DC motor can be changed by reversing either the field or armature connections but not both. This
is commonly done with a special set of contactors (direction contactors).The effective voltage
can be varied by inserting a series resistor or by an electronically controlled switching device
made of thyristors, transistors, or, formerly, mercury arc rectifiers.

28
 In a circuit known as a chopper, the average voltage applied to the motor is varied by
switching the supply voltage very rapidly. As the "on" to "off" ratio is varied to alter the average
applied voltage, the speed of the motor varies. The percentage "on" time multiplied by the supply
voltage gives the average voltage applied to the motor. Therefore, with a 100 V supply and a
25% "on" time, the average voltage at the motor will be 25 V. During the "off" time, the
armature's inductance causes the current to continue through a diode called a "flyback diode", in
parallel with the motor. At this point in the cycle, the supply current will be zero, and therefore
the average motor current will always be higher than the supply current unless the percentage
"on" time is 100%. At 100% "on" time, the supply and motor current are equal. The rapid
switching wastes less energy than series resistors. This method is also called pulse-width
modulation (PWM) and is often controlled by a microprocessor. An output filter is sometimes
installed to smooth the average voltage applied to the motor and reduce motor noise.

Since the series-wound DC motor develops its highest torque at low speed, it is often
used in traction applications such as electric locomotives, and trams. Another application is
starter motors for petrol and small diesel engines. Series motors must never be used in
applications where the drive can fail (such as belt drives). As the motor accelerates, the armature
(and hence field) current reduces. The reduction in field causes the motor to speed up until it
destroys itself. This can also be a problem with railway motors in the event of a loss of adhesion
since, unless quickly brought under control, the motors can reach speeds far higher than they
would do under normal circumstances. This can not only cause problems for the motors
themselves and the gears, but due to the differential speed between the rails and the wheels it can
also cause serious damage to the rails and wheel treads as they heat and cool rapidly. Field
weakening is used in some electronic controls to increase the top speed of an electric vehicle.

The simplest form uses a contactor and field-weakening resistor; the electronic control
monitors the motor current and switches the field weakening resistor into circuit when the motor
current reduces below a preset value (this will be when the motor is at its full design speed).
Once the resistor is in circuit, the motor will increase speed above its normal speed at its rated
voltage. When motor current increases, the control will disconnect the resistor and low speed
torque is made available.

29
 One interesting method of speed control of a DC motor is the Ward Leonard control. It is
a method of controlling a DC motor (usually a shunt or compound wound) and was developed as
a method of providing a speed-controlled motor from an AC supply, though it is not without its
advantages in DC schemes. The AC supply is used to drive an AC motor, usually an induction
motor that drives a DC generator or dynamo. The DC output from the armature is directly
connected to the armature of the DC motor (sometimes but not always of identical construction).
The shunt field windings of both DC machines are independently excited through variable
resistors. Extremely good speed control from standstill to full speed, and consistent torque, can
be obtained by varying the generator and/or motor field current. This method of control was the
de facto method from its development until it was superseded by solid state thyristor systems.

It found service in almost any environment where good speed control was required, from
passenger lifts through to large mine pit head winding gear and even industrial process
machinery and electric cranes. Its principal disadvantage was that three machines were required
to implement a scheme (five in very large installations, as the DC machines were often
duplicated and controlled by a tandem variable resistor). In many applications, the motor-
generator set was often left permanently running, to avoid the delays that would otherwise be
caused by starting it up as required. Although electronic (thyristor) controllers have replaced
most small to medium Ward-Leonard systems, some very large ones (thousands of horsepower)
remain in service. The field currents are much lower than the armature currents, allowing a
moderate sized thyristor unit to control a much larger motor than it could control directly. For
example, in one installation, a 300 amp thyristor unit controls the field of the generator. The
generator output current is in excess of 15,000 amperes, which would be prohibitively expensive
(and inefficient) to control directly with thyristors.

3.20 INFRARED SENSOR

In this IR detector and transmitter circuit the IC 555 is working under a stable mode. The
pin 4 i.e. reset pin is when grounded via IR receiver the pin 3 output is low.

As soon as the IR light beam transmitted is obstructed, a momentary pulse actuates the
relay output (or LED).The IR transmitter is simple series connected resistor network from
battery. The timing capacitor connected to pin 2 and ground can varied as per requirement

30
3.21 SWITCHED-MODE POWER SUPPLY

A switched-mode power supply (switching-mode power supply, SMPS, or switcher) is an

electronic power supply that incorporates a switching regulator to convert electrical power
efficiently. Like other power supplies, an SMPS transfer’s power from a source like the
electrical power grid to a load (such as a personal computer) while
converting voltage and current characteristics. An SMPS is usually employed to efficiently
provide a regulated output voltage, typically at a level different from the input voltage.

Unlike a linear power supply, the pass transistor of a switching mode supply continually
switches between low-dissipation, full-on and full-off states, and spends very little time in the
high dissipation transitions (which minimizes wasted energy). Ideally, a switched-mode power
supply dissipates no power

Switched mode power supply circuitry

Unlike a linear power supply, the pass transistor of a switching mode supply continually
switches between low-dissipation, full-on and full-off states, and spends very little time in the
high dissipation transitions (which minimizes wasted energy). Ideally, a switched-mode power
supply dissipates no power. Voltage regulation is achieved by varying the ratio of on-to-off time.

31
In contrast, a linear power supply regulates the output voltage by continually dissipating power
in the pass transistor. This higher power conversion efficiency is an important advantage of a
switched-mode power supply. Switched-mode power supplies may also be substantially smaller
and lighter than a linear supply due to the smaller transformer size and weight.

Switching regulators are used as replacements for the linear regulators when higher
efficiency, smaller size or lighter weights are required. They are, however, more complicated,
their switching currents can cause electrical noise problems if not carefully suppressed, and
simple designs may have a linear regulator provides the desired output voltage by dissipating
excess power in ohmic losses (e.g., in a resistor or in the collector–emitter region of a pass
transistor in its active mode). A linear regulator regulates either output voltage or current by
dissipating the excess electric power in the form of heat, and hence its maximum power
efficiency is voltage-out/voltage-in since the volt difference is wasted. In contrast, a switched-
mode power supply regulates either output voltage or current by switching ideal storage
elements, like inductors and capacitors, into and out of different electrical configurations. Ideal
switching elements (e.g., transistors operated outside of their active mode) have no resistance
when "closed" and carry no current when "open", and so the converters can theoretically operate
with 100% efficiency (i.e., all input power is delivered to the load; no power is wasted as
dissipated heat).

For example, if a DC source, an inductor, a switch, and the corresponding electrical

ground are placed in series and the switch is driven by a square wave, the peak-to-peak voltage
of the waveform measured across the switch can exceed the input voltage from the DC source.
This is because the inductor responds to changes in current by inducing its own voltage to
counter the change in current, and this voltage adds to the source voltage while the switch is
open.

If a diode-and-capacitor combination is placed in parallel to the switch, the peak voltage

can be stored in the capacitor, and the capacitor can be used as a DC source with an output
voltage greater than the DC voltage driving the circuit. This boost converter acts like a step-up
transformer for DC signals. A buck–boost converter works in a similar manner, but yields an
output voltage which is opposite in polarity to the input voltage.

Other buck circuits exist to boost the average output current with a reduction of voltage.

32
3.22 BATTERY

An electrochemical battery - or, more precisely, a "cell" - is a device in which the reaction
between two substances can be made to occur in such a way that some of the chemical energy is
converted to useful electricity. When the cell can only be used once, it is called a "primary" cell.
When the chemical reaction can be reversed repeatedly by applying electrical energy to the cell,
it is called a "secondary" cell and can be used in an accumulator or "storage" battery.

Certain cells are capable of only a few charge-discharge cycles and are, therefore, technically
"secondary" cells. Such is the case with certain silver oxide-zinc batteries. These batteries are not
capable of the repeated cycling required of a satellite battery system, and are, therefore,
considered to be "rechargeable primary" rather than storage batteries. To define a battery in
another way, it is an arrangement whereby an "electrochemical" reaction can be made to take
place so that the "electrical" part of the reaction proceeds via the metallic path of the external
circuit, while the "chemical" part of the reaction occurs via ionic conduction through electrolyte.

The type of chemical reaction that can be used in an electrochemical cell is known as an
"oxidation-reduction" reaction - a reaction in which one chemical species gives electrons to
another. By separating the two species and controlling the flow of ions between them, battery
engineers make devices in which essentially all of these electrons can be made to flow through
an external circuit, thereby converting most of the chemical energy to electrical energy during
the discharge of the cell.

Some of the components common to all cells are:

1. The "cathode" or "positive" electrode, which consists of a mass of "electron-receptive"

chemical held in intimate contact with a metallic "plate" through which the electrons arrive
from the external circuit.
2. The "anode" or "negative" electrode, which consists of another chemical which readily gives
up electrons - an "electron donor" - similarly held in close contact with a metallic member
through which electrons can be conducted to the external circuit.

33
3. The "electrolyte," usually a liquid solution that permits the transfer of mass necessary to the
overall reaction. This movement takes place by "migration" of "ions" - positively or
negatively charged molecular fragments - from anode to cathode and from cathode

BATTERY

3.23 ROLLER

Roller is made up of rubber material. If the roller will be connected to the dc motor and
the motor is rotated the roller also rotated. It is used to pull the paper from input tray to output
tray.

34
 CHAPTER-4
WORKING PRINCIPLE

The housing and paper counting mechanism from a Dell brand Lexmark printer was selected to
be utilized in the final design. It was chosen for 4 reasons:

1) The counting mechanism’s reliability and accuracy.

2) The large input paper tray which can hold up to 500 sheets of paper.

3) The ease of disassembly.

4) The professional appearance and layout of the device.

The professional appearance is important because it resembles office products the

students may encounter in the future such as industrial printers, fax machines, and copiers. The
original electronics were removed and space was cut in the interior housing in order to make
room for our electronics and controls.

A DC motor attachment was designed so that the paper-handling mechanism could be

run by a motor relocated in the lower housing. This motor attachment was adapted and press fit
onto the original drive shaft. The new shaft/motor assembly was relocated inside the housing
and was mounted with a specially designed bracket to minimize vibration and noise. The roller is
connected to the dc motor attachment and the motor is on the roller will be rotated.

The roller is fitted on the paper in input tray. Sensors were mounted in the end of the
input tray and on the housing to add functionality and safety features to the paper counter. An
optical/tactile counting sensor (an optical sensor which registers the movement of a tactile arm)
was mounted on the housing to detect and count paper as it passes from the lower input tray to
the upper output tray. A tactile ‘Empty Tray’ sensor was mounted in the input tray to stop the
motors automatically if the tray runs out of paper.

35
 A 6 Volt power source is used to run the sensors, motors, and display. The LCD display to
indicate the result to output tray.

36
 CHAPTER-5
ADVANTAGES, DISADVANTAGES & APPLICATIONS

5.1 ADVANTAGES
 Paper counting will be easy.
 Human work can be reduced.
 Time can be saved.
 Low cost.
 Time saving process

5.2 DISADVANTAGES
 The paper may be fold one on another
 For accurate counting, the second time process of paper was required.

5.3 APPLICATIONS
 Paper mill to package the paper
 School and colleges
 Stationeries
 Office purpose
 Paper recycling industries

CHAPTER-6
COST ESTIMATION

37
CHAPTER-7
CONCLUSION

38
 Every project work has a thought or purpose behind it. Our project may not promise to
form the best Machine but it certainly promises to be able to be used as the base for further
developments. The main feature of the project is its portability and adaptability. Since it is
implemented in small size this enables it to be portable and the ability to handle very easily any
kind of places. The machine implementation of the same can be used for many purposes like
reducing man work, in industries, in institutes etc. It will count thousands of pages easily by the
machine.

The project has an inbuilt capability to be operated in the automatic mode and hence the
need for automation can also be fulfilled. We implemented a prototype that’s why its size
somewhat large but in future its size and circuitry complexity will reduce. More features must be
added with advance circuitry. Here to get more accuracy different types of sensors are used. In
future the machine will become more digitalized to user reliable. It is compatible to new
mechanism.

REFERENCES

 P. Srisuresh and K. Egevang, “Traditional IP network address translator (traditional

NAT),” RFC 3022, Internet Engineering Task Force, Jan. 2001.

39
 T. Hain, “Architectural implications of NAT,” RFC 2993, Internet Engineering Task
Force, Nov. 2000.

 J. Postel, “Internet protocol,” RFC 791, Internet Engineering Task Force, Sept. 1981.

 Ratul Mahajan, Neil T. Spring, and David Wetherall, “Measuring ISP topologies with
Rocket fuel,” in Proceedings of SIGCOMM 2002, 2002, to appear.

 J. C. Mogul and S. E. Deering, “Path MTU discovery,” RFC 1191, Internet Engineering
Task Force, Nov. 1990.

 M. Holdrege and P. Srisuresh, “Protocol complications with the IP network address

translator,” RFC 3027, Internet Engineering Task Force, Jan. 2001.

 D. Senie, “Network address translator (NAT)-friendly application design guidelines,”

RFC 3235, Internet Engineering Task Force, Jan. 2002.

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