01.introduction To Mathematics For Mechanical Engineering PDF

3
Introduction
1. Introduction to Mathematics
Part A 1
for Mechanical Engineering
Ramin S. Esfandiari
1.1 Complex Analysis .................................. 4

This chapter is concerned with fundamental math-
1.1.1 Complex Numbers ........................ 4
ematical concepts and methods pertaining to
1.1.2 Complex Variables and Functions ... 7
mechanical engineering. The topics covered in-
clude complex analysis, differential equations, 1.2 Differential Equations ........................... 9
Laplace transformation, Fourier analysis, and lin- 1.2.1 First-Order Ordinary
ear algebra. These basic concepts essentially act as Differential Equations ................... 9
tools that facilitate the understanding of various 1.2.2 Numerical Solution of First-Order
ideas, and implementation of many techniques, Ordinary Differential Equations ...... 10
involved in different branches of mechanical en- 1.2.3 Second- and Higher-Order,
Ordinary Differential Equations ...... 11
gineering. Complex analysis, which refers to the
study of complex numbers, variables and func- 1.3 Laplace Transformation......................... 15
tions, plays an important role in a wide range of 1.3.1 Inverse Laplace Transform ............. 16
areas from frequency response to potential theory. 1.3.2 Special Functions ......................... 18
The significance of ordinary differential equations 1.3.3 Laplace Transform of Derivatives
(ODEs) is observed in situations involving the rate and Integrals ............................... 21
of change of a quantity with respect to another. 1.3.4 Inverse Laplace Transformation...... 22
A particular area that requires a thorough knowl- 1.3.5 Periodic Functions ........................ 23
edge of ODEs is the modeling, analysis, and control 1.4 Fourier Analysis .................................... 24
of dynamic systems. Partial differential equations 1.4.1 Fourier Series............................... 24
(PDEs) arise when dealing with quantities that 1.4.2 Fourier Transformation ................. 25
are functions of two or more variables; for in-
1.5 Linear Algebra...................................... 26
stance, equations of motions of beams and plates.
1.5.1 Vectors and Matrices..................... 27
Higher-order differential equations are generally 1.5.2 Eigenvalues and Eigenvectors ........ 30
difficult to solve. To that end, the Laplace transfor- 1.5.3 Numerical Solution
mation is used to transform the data from the time of Higher-Order Systems of ODEs .... 32
domain to the so-called s-domain, where equa-
tions are algebraic and hence easy to treat. The References .................................................. 33
solution of the differential equation is ultimately
obtained when information is transformed back transformation maps information from the time to
to time domain. Fourier analysis is comprised of the frequency domain, and its extension leads to
Fourier series and Fourier transformation. Fourier the Laplace transformation. Linear algebra refers
series are a specific trigonometric series represen- to the study of vectors and matrices, and plays
tation of a periodic signal, and frequently arise in a central role in the analysis of systems with large
areas such as system response analysis. Fourier numbers of degrees of freedom.
4 Part A Fundamentals of Mechanical Engineering
1.1 Complex Analysis

Part A 1.1
Complex numbers, variables and functions are the main x = Re(z) = −1 and y = Im(z) = 2. A complex num-
focus of this section. We will begin with complex number with zero real part is known as pure imaginary, e.g.,
bers, their representations, as well as properties. The z = 4i. Two complex numbers are said to be equal if
idea is then extended to complex variables and their and only if their respective real and imaginary parts
functions. are equal. Addition of complex numbers is performed
component-wise, that is, if z 1 = x1 + iy1 and z 2 = x2 +
1.1.1 Complex Numbers iy2 , then
A complex number z appears in the rectangular form z 1 + z 2 = (x1 + iy1 ) + (x2 + iy2 )
= (x1 + x2 ) + i(y1 + y2 ) . (1.2)
z = x + iy ,
√ Multiplication of two complex numbers is performed in
i = −1 = imaginary number , (1.1) the same way as two binomials with the provision that
i2 = −1, i3 = −i, i4 = 1, etc. need be taken into account,
where x and y are real numbers, called the real and that is,
imaginary parts of z, respectively, and denoted by
x = Re(z), y = Im(z). For example, if z = −1 + 2i, then z 1 z 2 = (x1 + iy1 )(x2 + iy2 )
= x1 x2 + iy1 x2 + ix1 y2 + i2 y1 y2
Imaginary axis = (x1 x2 − y1 y2 ) + i(x1 y2 + x2 y1 ) . (1.3)
Complex Plane
Since complex numbers consist of a real part and an
imaginary part, they have a two-dimensional character,
and hence may be represented geometrically as points
z = x + iy
y in a Cartesian coordinate system, known as the complex
plane. The x-axis of the complex plane is the real axis,
and its y-axis is called the imaginary axis, (Fig. 1.1).
Noting that z = x + iy is uniquely identified by an or-
dered pair (x, y) of real numbers, we can represent z as
a two-dimensional (2-D) vector in the complex plane,
0 x Real axis with initial point 0 and terminal point z = x + iy; in
other words, the position vector of the point z. The
Fig. 1.1 Geometrical representation of complex numbers –
imaginary number i, for instance, can be identified by
the complex plane
(0, 1). So, the concept of vector addition also applies to
the addition of complex numbers. For that, let us con-
Imaginary axis sider z 1 = −2 + 3i and z 2 = 3 + i in Fig. 1.2. It is then
evident that their sum, z 1 + z 2 = 1 + 4i, is exactly what
4i 1 + 4i
we would obtain by adding the corresponding position
vectors of z 1 and z 2 .
3i –2 + 3i The magnitude of a complex number z = x + iy is
defined as

2i
|z| = x 2 + y2 . (1.4)
i 3+i Geometrically, |z| is the distance from z to the origin of

the complex plane. If z is real, it must be located on the
x-axis, and its magnitude is equal to its absolute value.
0
–2 –1 0 1 2 3 Real axis If z is pure imaginary (z = iy), then it is on the y-axis,
and |z| = |y|. The quantity |z 1 − z 2 | gives the distance
Fig. 1.2 Addition of complex numbers by vector addition between z 1 and z 2 (Fig. 1.3).
Introduction to Mathematics for Mechanical Engineering 1.1 Complex Analysis 5
we conclude that
Part A 1.1
y
1
x = Re(z) = (z + z̄) ,
z2 2
1
y = Im(z) = (z − z̄) . (1.7)
z1 – z2
2i
Division of Complex Numbers

Let us consider z 1 /z 2 where z 1 = x1 + iy1 and z 2 = x2 +
iy2 (= 0). Multiply the numerator and the denominator
z1
by the conjugate of the denominator, that is, z̄ 2 = x2
−iy2 . Then, by (1.6), the resulting denominator is sim-
ply |z 2 |2 = x22 + y22 , a real number. In summary,
0 x
x1 + iy1 x1 + iy1 x2 − iy2
=
Fig. 1.3 Distance between two complex numbers x2 + iy2 x2 + iy2 x2 − iy2
(x1 x2 + y1 y2 ) + i (y1 x2 − y2 x1 )
=
Example 1.1: Distance x22 + y22
Given z 1 = 5 + 2i and z 2 = −1 + 10i, then x1 x2 + y1 y2 y1 x2 − y2 x
= + 2 i (1.8)
|z 1 − z 2 | = |(5 + 2i) − (−1 + 10i)| x22 + y22 x2 + y22

= |6 − 8i| = 62 + (−8)2 = 10 where the outcome is represented in the standard rect-
angular form.
Addition of complex numbers obeys the triangle in- Example 1.3: Division of complex numbers
equality, Perform the following division of complex numbers,
and express the result in the standard rectangular form:
|z 1 + z 2 | ≤ |z 1 | + |z 2 | . (1.5)
2−i
Given a complex number z = x + iy, its conjugate, de- −1 + 4i
noted by z̄, is defined as z̄ = x − iy. An immediate result
is that the product of a complex number (z = 0) and its Solution. Multiplication and division by the conjugate
conjugate is a positive, real number, equal to the square of the denominator, yields
of its magnitude, that is, (2 − i)(−1 − 4i) −6 − 7i
=
(−1 + 4i)(−1 − 4i) 17
z z̄ = (x + iy)(x − iy) = x 2 + y2 . (1.6)
6 7
= − −i
Geometrically, a complex number and its conjugate are 17 17
reflections of one another about the real axis; (Fig. 1.4).
Example 1.2: Conjugation y
Given z = −1 + 2i, we have
z = x + iy
z z̄ = (−1 + 2i)(−1 + 2i) = (−1 + 2i)(−1 − 2i)
= (−1)2 + (2)2 = 5 ,
√
which agrees with |z| = | − 1 + 2i| = 5. 0
Complex conjugation is extremely useful in com-

plex algebra. To begin with, noting that
z + z̄ = (x + iy) + (x − iy) = 2x z = x – iy
and x
z − z̄ = (x + iy) − (x − iy) = 2iy Fig. 1.4 A complex number and its conjugate
Polar Representation of Complex Numbers The angle θ is measured from the positive real axis and,
Part A 1.1
Although the standard rectangular form is suitable in by convention, is regarded as positive in the sense of
certain instances, it is quite inconvenient in most oth- the counterclockwise (ccw) direction. It is measured in
ers. For example, imagine the simplification of (−2 + radians (rad) and is determined in terms of integer mul-
3i)10 . Situations of this type require a special form tiples of 2π. The specific value of θ that lies in the
that simplifies the complex algebra. The polar form of interval (−π, π] is called the principal value of arg z
a complex number, as suggested by its name, uses the and is denoted by arg z. In engineering analysis, it is
polar coordinates to represent a complex number in the also common to express the polar form of z as
complex plane. Recall that any point in the plane can
be determined by a radial coordinate r and an angular z = r θ (1.14)
coordinate θ. So, the same holds for a complex number where denotes the angle.
z = x + iy = 0 in the complex plane, (Fig. 1.5). The rela-
tionship between the rectangular and polar coordinates Example 1.4: Phase via location
is given by Express z = 2
in polar form.
−1+i
x = r cos θ , y = r sin θ . (1.9)
We first introduce Euler’s formula, Solution. First, express z in standard rectangular form,
as
e = cos θ + i sin θ .
iθ
(1.10)
2 −1 − i −2 − 2i
Then, (1.9) and (1.10) yield z= = = −1 − i ,
−1 + i −1 − i 2
z = x + iy = r cos θ + i (r sin θ) = r eiθ indicating that z is located in the third quadrant of the
In summary, complex plane. Next, we use (1.13) to find

z = r eiθ , (1.11) −1 π
θ = tan−1 = 45◦ = rad .
which is called the polar form of the complex number z. −1 4
Here, the magnitude (or modulus) of z is defined by However, the only information this provides is
√ that the (smallest) angle between OA and the real
r = |z| = x 2 + y2 = z z̄ (1.12)
axis (Fig. 1.6) is 45◦ . Since z is in the third quad-
and the phase (or argument) of z is rant, its actual phase is then 180 + 45 = 225◦ (π +

Im(z) π/4 = 5π/4 rad) if measured in the ccw direction, or
θ = arg z = tan−1 −135◦ (−3π/4 rad) in the clockwise (cw) direction. So,
Re(z)
the polar form of z can be written as
−1 y
= tan . (1.13)
x √ √ 5π
z = −1 − i = 2 ei(5π/4) or z= 2 .
4
Imaginary axis
Multiplication and Division in Polar Form. As cited
earlier, polar form substantially reduces complex alge-
y
y z = x + iy
0
r 45°
y = r sin θ –135°
θ
x A
–i
Real axis z = –1– i
x = r cos θ
–1 0 1 x
Fig. 1.5 Relation between the rectangular and polar forms
of a complex number Fig. 1.6 Example 1.4
Introduction to Mathematics for Mechanical Engineering 1.1 Complex Analysis 7
bra, in particular, multiplication and division. Consider
Part A 1.1
y
two complex numbers z 1 = r1 eiθ1 and z 2 = r2 eiθ2 .
Subsequently, z = r e iθ
z 1 z 2 = r1r2 ei(θ1 +θ2 ) or r1 r2 (θ1 + θ2 ) . (1.15) r

θ
This means the magnitude and phase of the product z 1 z 2
are 0
|z 1 z 2 | = r1r2 = |z 1 ||z 2 | and –θ

r
arg(z 1 z 2 ) = θ1 + θ2 = arg(z 1 ) + arg(z 2 ) (1.16)
Similarly, for division of complex numbers, we have z = r e -iθ
z1 r1 r1 0 x
= ei(θ1 −θ2 ) or (θ1 − θ2 ) . (1.17)
z2 r2 r2
Fig. 1.7 A complex number and its conjugate in polar form
so that

z 1 r1 |z 1 |
= Roots of a Complex Number √
z r = |z | and In real calculus, if a is a real number then n a has a sin-
2 2 2
z1 gle value. On the contrary, given a complex number
arg = θ1 − θ2 = arg(z 1 ) − arg(z 2 ) (1.18) z = 0, and
z2 √ a positive integer n, then the nth root of z,
written n √z, is multivalued. In fact, there are n different
values of n z, corresponding to each value of z = 0. For
Complex Conjugation in Polar Form. Given the polar a known z = r eiθ , it can be shown that [1.1, 2]
form of a complex number, z = r eiθ , its conjugate is
obtained as √ √ θ + 2kπ θ + 2kπ
n
z = n r cos + i sin ,
n n
z̄ = x − iy = r cos θ − i(r sin θ)
k = 0, 1, · · · , n − 1 . (1.21)
= r(cos θ − i sin θ)
Geometrically, these n values are described as follows:
r e−iθ .
Euler’s formula
= r[cos(−θ) + i sin(−θ)] =
(1.19) 1. they all lie √
on a circle centered at the origin with
a radius of n r, and
This result makes sense geometrically, since a complex 2. they are the n vertices of an n-sided regular polygon.
number and its conjugate are reflections of one another
through the real axis. Hence, they are equidistant from
Example 1.5: Fourth
√ roots of unity
the origin, that is, |z| = |z̄| = r, and the phase of one
We are seeking 4 z, where z = 1. Noting that z = 1 is
is the negative of the phase of the other, i. e., arg(z) =
on the positive real axis, one unit from the origin, we
− arg(z̄); Fig. 1.7. The important property of complex
conclude that r = 1 and θ = 0, hence z = 1 = 1 ei(0) .
conjugation (1.6) can now be confirmed in polar form,
Following (1.21), we find the four roots to be 1, i, −1,
as
and −i; Fig.
√ 1.8. Note that all four roots lie on a circle
z z̄ = (r eiθ )(r e−iθ ) = r 2 = |z|2 . of radius 4 1 = 1 centered at the origin (the so-called
unit circle), and are the vertices of a regular four-sided
polygon, as asserted.
Integer Powers of a Complex Number
The effectiveness of the polar form may further be
demonstrated when raising a complex number to an 1.1.2 Complex Variables and Functions
integer power. Letting z = r eiθ , then
If x or y or both vary, then z = x + iy is referred to
z n = (r eiθ )n = r n einθ as a complex variable. The most well-known complex
Euler’s formula n
= r (cos nθ + i sin nθ) , (1.20)
variable is the Laplace variable (Sect. 1.3). Letting S be
a set of complex numbers, a function f defined on S
so that Re(z n ) = r n cos nθ and Im(z n ) = r n sin nθ. is a rule, which assigns a complex number w to each
neighborhood of z 0 . A function is analytic in a domain

Part A 1.1
y
if it is analytic at all points of that domain. Analytic
i functions arise in such areas as fluid flow and complex
i
potentials.
Test of Analyticity: Cauchy–Riemann Equations.

Suppose f (z) = u(x, y) + iv(x, y) is defined and con-
-1 1
0 tinuous in a neighborhood of some point z = x + iy,
and that f (z) exists at that point. Then, the first par-
tial derivatives of u and v with respect to x and y
(that is, u x , u y , vx , v y ) exist at that point and satisfy the
Cauchy–Riemann equations
-i
-i
u x = vy , vx = −u y . (1.23)
-1 0 1 x
Fig. 1.8 Locations of the fourth roots of unity Consequently, if f (z) is analytic in some domain R,
then the Cauchy–Riemann equations hold at every point
z ∈ S. The notation is w = f (z) and the set S is called of R.
the domain of definition of f . As an example, the do-
Example 1.6: Cauchy–Riemann equations
main of the function w = z/(3 − z) is any region that
Decide whether f (z) = (i − 2)z 2 − 2iz + i is analytic.
does not contain the point z = 3. Because z assumes dif-
ferent values from S, it is clearly a complex variable.
Solution. Inserting z = x + iy, we find
Since w is complex, it must have a real part u and an
imaginary part v, or w = u + iv. Also w = f (z) implies u(x, y) = −2xy − 2x 2 + 2y2 + 2y ,
that w is dependent on z = x + iy. Therefore, w depends
on x and y, which means u and v depend on x and y, or v(x, y) = x 2 − y2 − 4xy − 2x + 1 .
w = f (z) = u(x, y) + iv(x, y) . (1.22) Partial differentiation yields

In real calculus, much can be learned about a function u x = −2y − 4x , v y = −2y − 4x ,
through its graph. However, when z and w are both
vx = 2x − 4y − 2 , u y = −2x + 4y + 2 ,
complex, such a convenient graph of w = f (z) is no
longer available. This is because each z and w is lo- so that the Cauchy–Riemann equations hold for all z,
cated in a plane rather than on a line; more exactly, and f (z) is analytic for all z.
each z is in the xy-plane and each w in the uv-plane.
However, a function f can still be thought of as a map- Cauchy–Riemann Equations in Polar Form. In many
ping (or transformation) that defines correspondence cases it is advantageous to use the polar form of the
between points z = (x, y) and w = (u, v). Then, the im- Cauchy–Riemann equations to test the analyticity of
age of a point z ∈ S is the point w = f (z), and the set of a function. The idea is to express z in polar form
images of all points z ∈ S that are mapped by f is the
range of f , and denoted by z = r eiθ = r(cos θ + i sin θ)
{w|w = f (z), z ∈ S} . so that the function f (z) = u(x, y) + iv(x, y) can be
written as f (z) = u(r, θ) + iv(r, θ). In this event, the
Analytic Functions Cauchy–Riemann equations in polar form can be de-
A function that is differentiable only at a single point rived as [1.1]
is not of practical interest to us. What is of interest, 1 1
however, is a function that is differentiable at a point u r = vθ , vr = − u θ . (1.24)
and an entire neighborhood of that point. A neighbor- r r
hood of a point is an open circular disk centered at the This is particularly useful when dealing with z m for
point. A complex function f (z) is analytic (or holomor- m ≥ 3, making it much easier to work with than
phic) at a point z 0 if it is differentiable throughout some (x + iy)m .
Introduction to Mathematics for Mechanical Engineering 1.2 Differential Equations 9
1.2 Differential Equations
Part A 1.2
Mathematical models of dynamic systems – mechani- 1.2.1 First-Order Ordinary
cal, electrical, electromechanical, liquid-level, etc. – are Differential Equations
represented by differential equations [1.3]. Therefore,
it is imperative to have a thorough knowledge of their First-order ODEs generally appear in the implicit form
basic properties and solution techniques. In this sec-
tion we will discuss the fundamentals of differential F(x, y, y ) = 0 . (1.26)
equations, specifically, ordinary differential equations
(ODEs), and present analytical and numerical methods For example, y + y2 = cos x can be expressed in the
to solve them. Differential equations are divided into above form with F(x, y, y ) = y + y2 − cos x. In other
two general categories: ordinary differential equations cases, the equation may be written explicitly as
and partial differential equations (PDEs). An equation
involving an unknown function and one or more of its y = f (x, y) . (1.27)
derivatives is called a differential equation. When there
An example would be y + 2y = ex where f (x, y)
is only one independent variable, the equation is called
= ex − 2y. A function y = s(x) is a solution of the first-
an ordinary differential equation (ODE). For example,
order ODE in (1.26) on a specified (open) interval if
y + 2y = ex is an ODE involving the unknown func-
it has a derivative y = s (x) and satisfies (1.26) for all
tion y(x), its first derivative y = dy/ dx, as well as
values of x in the given interval. If the solution is in
a given function ex . Similarly, xy − yy = sin x is an
the form y = s(x), then it is called an explicit solu-
ODE relating y(x) and its first and second derivatives
tion. Otherwise, it is in the form S(x, y) = 0, which is
with respect to x, as well as the function sin x. While
known as an implicit solution. For example, y = 4 e−x/2
dealing with time-varying functions – as in many phys-
is an explicit solution of 2y + y = 0. It turns out that
ical applications – the independent variable x will be
a single formula y = k e−x/2 involving a constant k = 0
replaced by t, representing time. In that case, the rate
generates all solutions of this ODE. Such formula is
of change of the quantity y = y(t) with respect to the
referred to as a general solution, and the constant is
independent variable t is denoted by ẏ = dy/ dt. If the
known as the parameter. When a specific value is
unknown function is a function of more than one inde-
assigned to the parameter, a particular solution is ob-
pendent variable, e.g., u(x, y), the equation is referred
tained.
to as a partial differential equation. The derivative of the
highest order of the unknown function y(x) with respect
Initial-Value Problem (IVP)
to x is the order of the ODE; for instance, y + 2y = ex
A first-order initial-value problem (IVP) appears in the
is of order one and xy − yy = sin x is of order two.
form
Consider an nth-order ordinary differential equation in
the form y = f (x, y) , y(x0 ) = y0 , (1.28)

an y (n)
+ an−1 y (n−1)
+ · · · + a1 y + a0 y = g(x) ,
where y(x0 ) = y0 , is called the initial condition.
(1.25)
where y = y(x) and y(n) = dn y/ dx n . If all coefficients Example 1.8: IVP

a0 , a1 , · · · , an are either constants or functions of the Solve the initial-value problem
independent variable x, then the ODE is linear. Other-
wise, the ODE is nonlinear. Based on this, y + 2y = ex 2y + y = 0 , y(2) = 3 .
describes a linear ODE, while xy − yy = sin x is non-
linear.
Solution. As mentioned earlier, a general solution is
Example 1.7: Order and linearity y = k e−x/2 . Applying the initial condition, we obtain
Consider 3y − (2x + 1)y + y
= ex .
Since the deriva-
Solve for k
tive of the highest order is three, the ODE is third order. y(2) = k e−1 = 3 ⇒ k = 3e .
Comparison with (1.25) reveals that n = 3, and a3 = 3,
a2 = −(2x + 1), a1 = 0, a0 = 1, and g(x) = ex . Thus, the Therefore, the particular solution is y = 3 e · e−x/2
ODE is linear. = 3 e1−x/2 .
Separable First-Order Ordinary Differential f (x) ≡ 0, then the ODE is called homogeneous, other-
Part A 1.2
Equations wise it is called nonhomogeneous.

A first-order ODE is referred to as separable if it can be
written as Solution of Linear First-Order ODEs
The general solution of (1.31) can be expressed as [1.1,
f (y)y = g(x) . (1.29) 4]

Using y = dy/ dx in (1.29), we have
y(x) = e−h(x) eh(x) f (x) dx + c ,
dy
f (y) = g(x) ⇒ f (y) dy = g(x) dx . (1.30)
dx where h(x) = g(x) dx . (1.32)
Integrating the two sides of (1.30) separately, yields
Note that the constant of integration in the calculation
f (y) dy = g(x) dx + c , c = const. of h is omitted because c accounts for all constants.
Example 1.10: Linear first-order ODE
Example 1.9: Separable ODE Find the particular solution to the initial-value problem
Solve the initial-value problem ex y = y2 , y(0) = 1. 2 ẏ + y = 4 e2t , y(0) = 1.
Solution. The ODE is separable and treated as Solution. Noting that t is now the independent vari-
able, we first rewrite the ODE to agree with the form
dy of (1.31), as
ex = y2
dx
1 1 1
Provided that y = 0
⇒ dy = dx ẏ + y = 2 e2t
y 2 ex 2
1 so that
⇒ − = − e−x + c
y 1
g = , f = 2 e2t .
(c = const.) 2
Solve for y 1
⇒ y(x) = −x , With h = g(t) dt = 12 dt = 12 t, the general solution is
e −c given by (1.32),
which is the general solution to the original differential
−t/2
equation. The specific value of c is determined via the y(t) = e e · 2 e dt + c
t/2 2t
given initial condition, as
⎫Initial condition 4
= e−t/2 2 e5t/2 dt + c = e2t + c e−t/2 .
y(0) = 1 ⎪ ⎬ 5
1
⇒ =1⇒c=0.
⎪
1 ⎭ 1 − c Applying the initial condition, we find y(0) = 45 +
y(0) = 1−c By gen. solution c = 1 ⇒ c = 15 . The particular solution is y(t) = 45 e2t +
1 −t/2
Substitution into the general solution yields the particu- 5e .
lar solution y(x) = ex .
1.2.2 Numerical Solution of First-Order
Linear First-Order Ordinary Ordinary Differential Equations
Differential Equations
A differential equation that can be expressed in the form Recall that a first-order ODE can appear in an implicit
form F(x, y, y ) = 0 or an explicit form y = f (x, y).
y + g(x)y = f (x) , (1.31) We will consider the latter, and assume that it is subject
where g and f are given functions of x, is called a lin- to a prescribed initial condition, that is,
ear first-order ordinary ODE. This of course agrees with y = f (x, y) , y(x0 ) = y0 , x0 ≤ x ≤ x N . (1.33)
what was discussed in (1.25) with slight changes in
notation. If f (x) = 0 for every x in the interval under If finding a closed-form solution of (1.33) is difficult
consideration, that is, if f is identically zero, denoted or impossible, we resort to a numerical solution. What
this means is that we find approximate values for the Further inspection reveals that RK4 produces the exact
Part A 1.2
solution y(x) at several points values (at least to five-decimal place accuracy) of the
x1 = x0 + h , x2 = x0 + 2h · · · xn solution at the mesh points.
= x0 + nh , · · · , x N = x0 + Nh
1.2.3 Second- and Higher-Order, Ordinary
known as mesh points, where h is called the step size. Differential Equations
Note that the mesh points are equally spaced. Among
many numerical methods to solve (1.33), the fourth- The application of basic laws such as Newton’s second
order Runge–Kutta method is most commonly used in law and Kirchhoff’s voltage law (KVL) leads to math-
practice. The difference equation for the fourth-order ematical models that are described by second-order
Runge–Kutta method (RK4) is derived as [1.5, 6] ODEs [1.3]. Although it is quite possible that the sys-
1 tem models contain nonlinear elements, in this section
yn+1 = yn + (q1 + 2q2 + 2q3 + q4 ) , (1.34) we will mainly focus on linear second-order differen-
6
n = 0, 1, · · · , N − 1 , tial equations. Nonlinear systems may be treated via
numerical techniques such as the fourth-order Runge–
where Kutta method (Sect. 1.2), or via linearization [1.3]. In
q1 = h f (xn , yn ) agreement with (1.25), a second-order ODE is said to
be linear if it can be expressed in the form
h q1
q2 = h f xn + , yn + ,
2 2 y + g(x)y + h(x)y = f (x) , (1.35)

h q2 where f , g, and h are given functions of x. Otherwise,
q3 = h f xn + , yn + ,
2 2 it is called nonlinear.
q4 = h f (xn + h, yn + q3 ) .
Homogeneous Linear Second-Order ODEs
If y1 and y2 are two solutions of the homogeneous linear
Example 1.11: Fourth-order Runge–Kutta method
ODE
Apply RK4 with step size h = 0.1 to solve y + y= 2x 2 ,
y(0) = 3, 0 ≤ x ≤ 1. y + g(x)y + h(x)y = 0 (1.36)
on some open interval, their linear combination
Solution. Knowing that f (xn , yn ) = −yn + 2xn2 , the y = c1 y1 + c2 y2 (c1 , c2 constants) is also a solution on
four function evaluations/step of the RK4 are the same interval. This is known as the principle of

q1 = h − yn + 2xn2 , superposition.

1 1 2 General Solution of Linear Second-Order
q2 = h − yn + q1 + 2 xn + h ,
2 2 ODEs – Linear Independence

1 1 2 A general solution of (1.36) is based on the idea of lin-
q3 = h − yn + q2 + 2 xn + h , ear independence of functions, which involves what is
2 2
known as the Wronskian. We first mention that a 2 × 2
q4 = h − (yn + q3 ) + 2(xn + h)2 . determinant (Sect. 1.5.1) is evaluated as

Upon completion of each step, yn+1 is calculated p q
= ps − qr .
by (1.34). So, we start with n = 0, corresponding to r s
x0 = 0 and y0 = 3, and continue the process up to
If each of the functions y1 (x) and y2 (x) has at least
n = 10. Numerical results are generated as
a first derivative, then their Wronskian is denoted by
y(0) = 3 , W(y1 , y2 ) and is defined as the 2 × 2 determinant

y(0.1) = 2.7152 ,
y1 y2
y(0.2) = 2.4613 , W(y1 , y2 ) = = y1 y2 − y2 y1 . (1.37)
y1 y2
y(0.3) = 2.2392 , · · · ,
If there exists a point x ∗ ∈ (a, b) where W = 0, then y1
y(0.9) = 1.6134 , and y2 are linearly independent on the entire interval
y(1) = 1.6321 . (a, b).
Example 1.12: Independent solutions – the Wronskian Since eλx = 0 for any finite values of x and λ, then
Part A 1.2
The functions y1 = e2x and y2 = e−3x are linearly in-

dependent for all x because their Wronskian is
λ 2 + a1 λ + a2 = 0

y1 y2 e2x e−3x
W(y1 , y2 ) = = 2x λ1 = 1
− a1 + a12 − 4a2
y1 y2 2 e −3 e−3x Solve the
2
⇒ .
= −5 e−x = 0 characteristic equation
λ2 = 1
2 − a1 − a12 − 4a2
for all x.
(1.40)
If y1 and y2 are two linearly independent solutions
of (1.36) on the interval (a, b), they form a basis of so- The solutions λ1 and λ2 of the characteristic equa-
lutions for (1.36) on (a, b). A general solution of (1.36) tion are the characteristic values. The assumption was
on (a, b) is a linear combination of the basis elements, y = eλx , hence the solutions of (1.39) are y1 = eλ1 x and
that is, y2 = eλ2 x . To find a general solution of (1.39), the two
independent solutions must be identified. But this de-
y = c1 y1 + c2 y2 (c1 , c2 constants) . (1.38) pends on the nature of the characteristic values λ1 and
λ2 , as discussed below.
Example 1.13: General solution, basis ` ´
It can be easily verified that y1 = e2x and y2 = e−3x Case 1: Two Distinct Real Roots a21 − 4a2 > 0‚λ1 = λ2 .
are solutions of y + y − 6y = 0 for all x. They are also In this case, the solutions y1 = eλ1 x and y2 = eλ2 x are
linearly independent by Example 1.12. Consequently, linearly independent, as may easily be verified. Thus,
y1 = e2x and y2 = e−3x form a basis of solutions for they form a basis of solution for (1.39). Therefore,
the ODE at hand, and a general solution for this ODE is a general solution is
y = c1 e2x + c2 e−3x (c1 , c2 constants).
y(x) = c1 eλ1 x + c2 eλ2 x . (1.41)
General solution — λ1 =λ2 , real
Example 1.14: Unique solution of an IVP
Find the particular solution of y + y − 6y = 0,
y(0) = −1, y (0) = 8. `
Case 2:´ Double (Real) Root a21 − 4a2 = 0‚ λ1 = λ2
= − 21 a1 . It can be shown [1.1] that the two lin-
Solution. By Example 1.13, a general solution is y =
early independent solutions are y1 = e−a1 x/2 and
c1 e2x + c2 e−3x . Differentiating and applying the initial
y2 = x e−a1 x/2 . Therefore,
conditions, we have
1 1
y(0) = c1 + c2 = −1 Solve the system c1 = 1 y(x) = c1 e− 2 a1 x + c2 x e− 2 a1 x
⇒ .
y (0) = 2c1 − 3c2 = 8
1
c2 = −2 = (c1 + c2 x) e− 2 a1 x . (1.42)
General solution — λ1 =λ2 , real
Therefore, the unique solution of the IVP is obtained
as y = e2x − 2 e−3x .
` ´
Case 3: Complex Conjugate Pair a21 − 4a2 < 0‚λ1 = λ̄2 .
characteristic values are given as λ1,2 = 2 (−a1 ±
1
Homogeneous Second-Order Differential The
Equations with Constant Coefficients a12 − 4a2 ). Since a12 − 4a2 < 0, we write
Consider a homogeneous linear second-order ODE with

constant coefficients, 1
λ1,2 = −a1 ± − 4a2 − a12
y + a1 y + a2 y = 0 (a1 , a2 constants) (1.39) 2

1 √
and assume that its solution is in the form y = eλx , = −a1 ± −1 4a2 − a12
where λ, known as the characteristic value, is to be 2
determined. Substitution into (1.39), yields
1
= −a1 ± i 4a2 − a1 = −σ ± iω ,
2
λ2 eλx + a1 λ eλx + a2 eλx = 0 2
√
⇒ eλx λ2 + a1 λ + a2 = 0 . (i = −1)
where Nonhomogeneous Linear Second-Order ODEs
Part A 1.2
1 Nonhomogeneous second-order ODEs appear in the
σ = a1 , form
2
1
ω= 4a2 − a12 . (1.43) y + g(x)y + h(x)y = f (x) ,
2
The two independent solutions are y1 = e−σ x cos ωx f (x) ≡ 0 . (1.46)
and y2 = e−σ x sin ωx, and a general solution of (1.39) A general solution for this equation is then obtained as
is obtained as
y(x) = yh (x) + yp (x) . (1.47)
y(x) = e−σ x (c1 cos ωx + c2 sin ωx) . Homogeneous solution Particular solution
General solution — λ1 =λ̄2 , complex conjugates
(1.44)
Homogeneous Solution yh (x). yh (x) is a general so-
Example 1.15: Case (3) lution of the homogeneous equation (1.36), and as
Solve y + 2y + 2y = 0, y(0) = 1, y (0) = 0. previously discussed, it is given by
Solution. We first find the characteristic equation and yh = c1 y1 + c2 y2 , (c1 , c2 constants)

the corresponding characteristic values, as
where y1 and y2 are linearly independent and form a ba-
λ2 + 2λ + 2 = 0 sis of solutions for (1.36). Note that the homogeneous
solution involves two arbitrary constants.
⇒λ1,2 = −1 ± i .
Complex conjugate pair, Case (3) Particular Solution yp (x). yp (x) is a particular solution
By (1.43), we identify σ = 1 and ω = 1, so that the of (1.46), and does not involve any arbitrary constants.
general solution by (1.44) is The nature of yp (x) depends on the nature of f (x), as
well as its relation to the independent solutions y1 and
y(x) = e−x (c1 cos x + c2 sin x) . y2 of the homogeneous equation.
Next, we differentiate this to obtain
Method of Undetermined Coefficients
y (x) = − e−x (c1 cos x + c2 sin x) When (1.46) happens to have constant coefficients and
+ e−x (−c1 sin x + c2 cos x) the function f (x) is of a special type – polynomial, ex-
ponential, sine and/or cosine or a combination of them
Finally, by the initial conditions, – then the particular solution can be obtained by the
method of undetermined coefficients as follows. Con-
y(0) = c1 = 1 c1 = 1
⇒ sider
y (0) = −c1 + c2 =0 c2 = 1
y + a1 y + a2 y = f (x) (a1 , a2 constants) . (1.48)
and the solution is y(x) = e−x (cos x + sin x).
Since the coefficients are constants, the homogeneous
Boundary-Value Problems (BVP). In certain appli- solution yh is found as before. So all we need to do is
cations involving second-order differential equations, to find the particular solution yp . We will make a proper
a pair of information is provided at the boundary points selection for yp based on the nature of f (x) and with the
of an open interval (a, b) on which the ODE is to be
solved. This pair is referred to as the boundary condi- Table 1.1 Selection of particular solution – the method of
tions, and the problem undetermined coefficients
Term in f (x) Proper choice of yp
y + a1 y + a2 y = 0 ,
an x n + . . . + a1 x + a0 Kn xn + . . . + K1 x + K0
y(a) = A , y(b) = B (1.45)
A eax K eax
Boundary conditions A sin ωx or A cos ωx K 1 cos ωx + K 2 sin ωx
is called a boundary-value problem (BVP). A eσ x sin ωx or A eσ x cos ωx eσ x (K 1 cos ωx + K 2 sin ωx)
aid of Table 1.1. This choice involves unknown coeffi- a homogeneous solution associated with a double root.
Part A 1.2
cients, which will be determined by substituting yp and Therefore, by special case II the modified choice is
its derivatives into (1.48). The details, as well as special Kx 2 e−x . Consequently, the particular solution is in the
cases that may occur, are given below. form
yp (x) = K 1 x + K 0 + K x 2 e−x .
Procedure. First term Second term
Step 1: Homogeneous Solution yh (x). Solve the homo- Substitution of yp and its derivatives into the nonhomo-
geneous equation y + a1 y + a2 y = 0 to find the two geneous ODE, and collecting terms, results in
independent solutions y1 and y2 , and the general solu-
tion yh (x) = c1 y1 (x) + c2 y2 (x). 2K e−x + K 1 x + K 0 + 2K 1 = x + 1 + 3 e−x .
Equating the coefficients of like terms, we have
Step 2: Particular Solution yp (x). For each term in f (x)
2K = 3 K = 32
choose a proper yp as suggested by Table 1.1. For in-
stance, if f (x) = x + 2 ex then pick yp = K 1 x + K 2 + K1 = 1 ⇒ K1 = 1
K ex . Note that, if instead of x we had 3x − 2, for ex- K 0 + 2K 1 = 1 K 0 = −1
ample, the choice of yp would still be the same because 3
they both represent first-degree polynomials. We then ⇒ yp (x) = x − 1 + x 2 e−x .
2
substitute our choice of yp , along with its derivatives,
into the original ODE to find the undetermined coeffi-
cients. Step 3: General Solution. The general solution is then
found as
3
Special cases. y(x) = (c1 + c2 x) e−x + x − 1 + x 2 e−x .
I. Suppose a term in our choice of yp coincides with 2
a solution (y1 or y2 ) of the homogeneous equation,
Step 4: Initial Conditions. Applying the initial condi-
and that this solution is associated with a simple
tions, we obtain c1 = 2 and c2 = 1. Finally, the solution
(i. e., nonrepeated) characteristic value. Then, make
to the IVP is
the modification by multiplying yp by x.
3
II. If a term in the choice of yp coincides with a so- y(x) = (2 + x) e−x + x − 1 + x 2 e−x .
lution of the homogeneous equation, and that this 2
solution is associated with a repeated characteristic
value, modify by multiplying yp by x 2 . Higher-Order Ordinary Differential Equations
Many of the techniques for the treatment of differential
Example 1.16: Special case II equations of order three or higher are merely extensions
Solve of those applied to second-order equations. Here we will
only discuss nth-order, linear nonhomogeneous ODEs
y + 2y + y = x + 1 + 3 e−x , y(0) = 1 , with constant coefficients, that is,
y (0) = 0 . y(n) + an−1 y(n−1) + · · · + a1 y + a0 y = f (x) ,
(1.49)
Step 1: Homogeneous Solution. The characteristic
equation (λ + 1)2 = 0 yields a double root λ = −1. This where a0 , a1 , · · · , an−1 are constants. As in the case
means y1 = e−x and y2 = x e−x , so that the homoge- of second-order ODEs, a general solution consists of
neous solution is yh (x) = (c1 + c2 x) e−x . the homogeneous solution and the particular solution.
For cases when f (x) is of a special type, the particu-
Step 2: Particular Solution. The right-hand side of the lar solution is obtained via the method of undetermined
ODE consists of two functions, coefficients.
x +1 and e−x . Method of Undetermined Coefficients. The idea intro-

First-degree polynomial
duced for second-order ODEs is now extended to find
The first term, x + 1, does not coincide with either y1 yp for (1.49). As before, a proper choice of yp is made
or y2 , so the proper choice by Table 1.1 is K 1 x + K 0 . assuming that f (x) consists of terms that are listed in
The second term involves e−x , which happens to be Table 1.1. If none of the terms in f (x) happens to be
Introduction to Mathematics for Mechanical Engineering 1.3 Laplace Transformation 15
an independent homogeneous solution, then no modi- Step 2: Particular Solution. Noting that f (x) = 1 + 12x
Part A 1.3
fication is necessary. Otherwise, the following special is a first-degree polynomial, we pick yp = K 1 x + K 0 .
cases need be taken into account. But x happens to be a homogeneous solution associated
with a double root (λ = 0). Hence, the modification is
Special Cases. yp = (K 1 x + K 0 )x 2 . Substituting this and its derivatives
1. If a term in our choice of yp coincides with a homo- into the original ODE, and simplifying, we arrive at
geneous solution, which corresponds to a simple
(6K 1 − 8K 0 ) − 24K 1 x = 1 + 12x
(nonrepeated) characteristic value, then we make
the modification by multiplying yp by x. 6K 1 − 8K 0 = 1 K = − 12
⇒ ⇒ 1
2. If a term in yp coincides with a solution of the −24K 1 = 12 K 0 = − 12
homogeneous equation, and this solution is associ- 1
ated with a characteristic value of multiplicity m, we ⇒ yp = − (x + 1)x 2
modify by multiplying yp by x m . 2
Example 1.17: Special case II Step 3: General Solution. Combination of yh and yp

Solve gives a general solution y = c1 + c2 x + c3 e4x − 12 (x +
1)x 2 .
y − 4y = 1 + 12x ,
y(0) = 0 , y (0) = 4 , y (0) = 15 . Step 4: Initial Conditions. Applying the initial con-
ditions to the general solution and its derivatives, we
obtain
Solution.
Step 1: Homogeneous Solution. Characteristic equa- y(0) = c1 + c3 = 0 c1 = −1
tion:
y (0) = c2 + 4c3 = 4 ⇒ c2 = 0
λ3 − 4λ2 = λ2 (λ − 4) = 0 ⇒ λ = 0, 0, 4 . y (0) = 16c3 − 1 = 15 c3 = 1
1 1
⇒ y(x) = −1 + e4x − x 3 − x 2 .
Therefore yh = c1 + c2 x + c3 e4x . 2 2
1.3 Laplace Transformation

In Sect. 1.2 we mainly learned to solve linear time- able, denoted by s. If a function f (t) is defined for all
invariant (LTI) ODEs without ever leaving the time t ≥ 0, then its Laplace transform is defined by
domain. In this section we introduce a systematic ap- Notation
proach to solve such ODEs in a more-expedient manner. F(s) = L[ f (t)]
The primary advantage gained here is that the arbitrary ∞
e−st f (t) dt
Definition
constants in the general solution need not be found sep- = (1.50)
arately. The idea is simple: in order to solve an ODE and
0
corresponding initial-value problem (IVP) or boundary-
value problem (BVP), transform the problem to the provided that the integral exists. The complex vari-
so-called s domain, in which the transformed problem is able s is the Laplace variable, and L is the Laplace
an algebraic one. This algebraic problem is then treated transform operator. It is common practice to denote
properly, and the data is ultimately transformed back to a time-dependent function by a lower-case letter, say,
time domain to find the solution of the original problem. f (t), and its Laplace transform by the same letter in
The transform function is a function of a complex vari- upper case, F(s).
1.3.1 Inverse Laplace Transform and a1 and a2 are constant scalars, then
Part A 1.3
L[a1 f 1 (t) + a2 f 2 (t)]

Suppose we are seeking the solution x(t) of an ODE.
The ODE is first transformed into the s domain by ∞
means of the operator L. In this domain, the transformed = e−st [a1 f 1 (t) + a2 f 2 (t)] dt
version of the ODE is an algebraic equation involving 0
the transform function X(s) of x(t). This equation is ∞ ∞
−st
then manipulated to find X(s), which in turn will be = a1 e f 1 (t) dt + a2 e−st f 2 (t) dt
transformed back into time domain to determine x(t).
0 0
This is done through the inverse Laplace transforma-
tion, as in Fig. 1.9. = a1 L[ f 1 (t)] + a2 L[ f 2 (t)]
= a1 F1 (s) + a2 F2 (s) . (1.52)
x(t) = L−1 [X(s)]
Consistent with (1.50), we have To establish the linearity of L−1 , take the inverse
Laplace transforms of the expressions on the far left and
f (t) = L−1 [F(s)] . (1.51) far right of (1.52) to obtain
a1 f 1 (t) + a2 f 2 (t) = L−1 [a1 F1 (s) + a2 F2 (s)] .
Example 1.18: Laplace transform
Given g(t) = 1 for t ≥ 0, find L[g(t)]. Noting that f 1 (t) = L−1 [F1 (s)] and f 2 (t) = L−1 [F2 (s)],
the result follows.
Solution
Following the definition given by (1.50), we have Example 1.20: Linearity of L
Find L[2 − 3 e4t ].
∞
L[g(t)] = L(1) = e−st dt Solution. Using the linearity of L, we write L[2 −
0 3 e4t ] = 2L[1] − 3L[ e4t ]. But, by Example 1.18 we

−st ∞ have L[1] = 1/s. And by Example 1.18 (with a = −4)
e 1
= = we have L[ e4t ] = 1/(s − 4). Thus L[2 − 3 e4t ] = 2/s −
−s t=0 s 3/(s − 4) = (−(s + 8))/(s(s − 4)).
for s > 0. Table of Laplace Transform Pairs. Laplace transforms
of several functions are listed in Table 1.2 at the end of
Example 1.19: Laplace transform
this section. We will refer to this frequently. For a better
Suppose h(t) = e−at (a = const) for t ≥ 0. Determine
understanding of the concepts, however, we try to derive
H(s).
the most fundamental results on our own.
Solution
Theorem 1.1: Shift on the s-axis. Suppose that F(s) =
By definition,
L[ f (t)] and that a is a constant. Then,
∞
H(s) = L[ e−at ] = e−st e−at dt L[ e−at f (t)] = F(s + a) . (1.53)
0
∞ ∞ Laplace
e−(s+a)t
= e−(s+a)t dt = transformation
−(s + a) t=0 Time domain
X(s) = L [x(t)]
s domain
0
Algebraic
1 ODE in x(t)
= equation
s+a in X(s)
for s + a > 0. x(t) = L–1 [X(s)]
Inverse Laplace
Linearity of Laplace transformation
and Inverse Laplace Transforms
The Laplace transform operator L is linear, that is, if the Fig. 1.9 Operations involved in the Laplace transformation
Laplace transforms of functions f 1 (t) and f 2 (t) exist, method
Table 1.2 Laplace transform pairs Table 1.2 Laplace transform pairs, continued
Part A 1.3
No. f (t) F(s) No. f (t) F(s)
√
1 Unit impulse δ(t) 1 1
π
32 1
e−at + 2 e 2 at sin 2 at − 6
3 1
2 1, unit step u s (t) 1/s 3a2 s 3 −a3
√
3 t, unit ramp u r (t) 1/s2 1
π
33 1
− e−at + 2 e 2 at sin 2 at + 6
3 s
4 δ(t − a) e−as
3a s 3 −a3
√
5 u(t − a) e−as /s 1
34 1
eat − 2 e− 2 at sin 3 π
2 at + 6
1
6 t n−1 , n = 1, 2, . . . (n − 1)!/sn 3a2 s 3 −a3
√
7 t a−1 , a > 0 Γ (a)/sa 1
35 1
e−at + 2 e− 2 at sin 3 π
2 at − 6
s
3a s 3 −a3
8 e−at 1
s+a
36 1
(cosh at sin at − sinh at cos at) 1
9 t e−at 1 4a3 s 4 +4a4
(s+a)2 1 s
37 sinh at sin at
10 t n e−at , n = 1, 2, . . . n! 2a2 s 4 +4a4
(s+a)n+1
−at
38 1
(sinh at − sin at) 1
11 1
b−a ( e − e−bt ) , a = b 1
(s+a)(s+b)
2a3 s 4 −a4
39 1
(cosh at − cos at) s
12 1 −at − b e−bt ) , a = b s 2a2 s 4 −a4
a−b (a e (s+a)(s+b)

13 1
ab 1 + a−b
1
(b e−at − a e−bt ) 1
s(s+a)(s+b) See Fig. 1.10. Alternatively, in terms of the inverse
14 1
(−1 + at + e−at ) 1 Laplace transform,
a2 s 2 (s+a)
15 1
(1 − e−at − at e−at ) 1 L−1 [F(s + a)] = e−at f (t) . (1.54)
a2 s(s+a)2
16 sin ωt ω Example 1.21: Shift on the s-axis

s 2 +ω2
Find L[ e3t cos t].
17 cos ωt s
s 2 +ω2
18 e−σt sin ωt ω Solution. Let f (t) = cos t so that F(s) = s/(s 2 + 1); see
(s+σ)2 +ω2 Table 1.2. Then, by (1.53) with a = −3,
ω2
19 e−σt cos ωt f (t)=cos t s−3
s(s 2 +ω2 )
L[ e3t cos t] = F(s − 3) =
20 1 − cos ωt s+σ a=−3 (s − 3)2 + 1
(s+σ)2 +ω2
ω3
21 ωt − sin ωt Differentiation and Integration
s 2 (s 2 +ω2 )
22 t cos ωt s 2 −ω2 of Laplace Transforms

(s 2 +ω2 )2
We now turn our attention to two specific types of
23 1
2ω t sin ωt s
situations: (1) L[t f (t)], (2) L[ f (t)/t]. In both cases,
(s 2 +ω2 )2
24 1
(sin ωt − ωt cos ωt) 1
2ω3 (s 2 +ω2 )2
s2 F (s)
2ω (sin ωt + ωt cos ωt)
1
25
(s 2 +ω2 )2
(Assuming a > 0)
26 1 1
ω2 sin ω2 t − ω1 sin ω1 t , 1

ω22 −ω21 1 s 2 +ω21 s 2 +ω22 a
ω21 = ω22
27 1
(cos ω1 t − cos ω2 t) , s

ω22 −ω21 s 2 +ω21 s 2 +ω22
ω21 = ω22 F (s + a) F (s)

a
28 sinh at
s 2 −a2
s
29 cosh at
s 2 −a2
1
30 1
sinh at − 1b sinh bt , 1
a2 −b2 a (s 2 −a2 )(s 2 −b2 )
a = b 0
s
31 1
[cosh at − cosh bt] , a = b s
a2 −b2 (s 2 −a2 )(s 2 −b2 )
Fig. 1.10 Shift on the s-axis (Theorem 1.1)
we assume that f (t) is such that F(s) = L[ f (t)] is Example 1.23: Theorem 1.3
Part A 1.3
either known directly from Table 1.2 or can be de- Show that
termined by other means. Either way, once F(s) is
sin ωt s
available, the two transforms labeled (1) and (2) will L = cot−1 .
be obtained in terms of the derivative and integral of t ω
F(s), respectively. Before presenting two key results
pertaining to these situations we make the follow- Solution. Comparing with (1.58), f (t) = sin ωt so that
ing definition. If a transform function is in the form F(s) = ω/(s2 + ω2 ). Subsequently,
F(s) = N(s)/D(s), then each value of s for which
D(s) = 0 is called a pole of F(s). A pole with a mul- ∞
sin ωt ω
tiplicity (number of occurrences) of one is known as L = dσ
a simple pole. t σ 2 + ω2
s
∞
Theorem 1.2: Differentiation of Laplace Transforms. 1 dσ
=
If L[ f (t)] = F(s) exists, then at any point except at the 1 + (σ /ω)2 ω
s
poles of F(s), we have σ ∞
d = tan−1
L[t f (t)] = − F(s) = −F (s) (1.55) ω σ=s
ds π s s
= − tan−1 = cot−1 .
or alternatively, 2 ω ω
t f (t) = −L−1 [F (s)] . (1.56)

1.3.2 Special Functions
The general form of (1.55) for n = 1, 2, 3, · · · is given
by Much can be learned about the characteristics of
dn a system based on its response to specific external
L[t n f (t)] = (−1)n F(s) = (−1)n F (n) (s) . (1.57) disturbances. To perform the response analysis, these
dsn disturbances must first be mathematically modeled,
Example 1.22: Differentiation of F(s) which is where special functions play an important role.
Find L[t sin 3t]. In this section we will introduce the step, ramp, pulse,
and impulse functions, as well as their Laplace trans-
Solution. Comparing with the left side of (1.55), we forms.
have f (t) = sin 3t so that F(s) = 3/(s2 + 9). Therefore,
Unit Step u(t)
d 3 The unit-step function (Fig. 1.11) is analytically defined
L[t sin 3t] = −
ds s2 + 9 as
6s ⎧
= 2 ⎪ if t > 0
(s + 9)2 ⎪
⎨1
u(t) = 0 if t < 0 . (1.60)
⎪
⎪
⎩
Theorem 1.3: Integration of Laplace transforms. If undefined (finite) if t = 0
L[ f (t)/t] exists, and the order of integration can be This may be physically realized as a constant signal (of
interchanged, then magnitude 1) suddenly applied to the system at time t =
∞ 0. By the definition of the Laplace transform, we find
f (t)
L = F(σ) dσ . (1.58)
t ∞
Notation
s
L[u(t)] = U(s) = e−st u(t) dt
Alternatively, 0
⎡ ⎤
∞ ∞
1
f (t) = tL−1 ⎣ F(σ) dσ ⎦ . (1.59) = e−st dt = . (1.61)
s
s 0
Theorem 1.4: Shift on the t-axis. Given that F(s) =
Part A 1.3
u (t)
L[ f (t)] exists, then
L[ f (t − a)u(t − a)] = e−as F(s) , (1.63)
1 or, alternatively,
L−1 [ e−as F(s)] = f (t − a)u(t − a) . (1.64)
Finding L[u(t − a)] via Theorem 1.4. We now have the

0 tools to determine L[u(t − a)]. In particular, comparing
L[u(t − a)] with the left-hand side of (1.63), we deduce
that f (t − a) = 1. Which implies that f (t) = 1, hence
0 t F(s) = 1/s. As a result,
Fig. 1.11 The unit step function u(t) e−as
L[u(t − a)] = . (1.65)
s
When the magnitude is some A = 1, we refer to the sig-
nal as a step function, denoted by Au(t). In this case, Unit Ramp ur (t). The unit ramp function (Fig. 1.13) is
analytically defined as
∞ ⎧
A ⎨t if t ≥ 0
L[Au(t)] = e−st A dt = .
s u r (t) =
0 ⎩0 if t < 0 .
When the unit step function occurs at some time a = Physically, this models a signal that changes linearly
0 (Fig. 1.12), it is denoted by u(t − a), and with a unit rate. By (1.50),
⎧ ∞
⎪
⎪ if t > a
⎨1 L[u r (t)]
Notation
= Ur (s) = t e−st dt
u(t − a) = 0 if t < a . (1.62)
⎪
⎪ 0
⎩ ∞ ∞ −st
undefined (finite) if t = a
e−st e
= t − dt
As before, if the magnitude happens to be A = 1, the −s t=0 −s
0
notation is modified to Au(t − a). To find the Laplace −st ∞
transform of u(t − a), we first need to discuss the shift e 1
= = 2. (1.66)
on the t-axis, see Theorem 1.4 below. −s t=0 s
2
u(t - a) ur (t)
1
1
1
0 0
0 a t 0 t
Fig. 1.12 The unit -step function occurring at t = a Fig. 1.13 The unit ramp function u r (t)
Note that u r (t) = tu(t). When the rate is A = 1, the sig- If the area is A = 1, the signal is called a pulse, written
Part A 1.3
nal is called a ramp function, denoted by Au r (t). In that Au p (t), and

case,
A(1 − e−st1 )
A
L[Au r (t)] = 2 . L[Au p (t)] = .
s st1
Unit Pulse up (t). The unit pulse function (Fig. 1.14) is Unit Impulse (Dirac Delta) δ(t). Consider the unit
defined as ⎧ pulse of Fig. 1.14 and let t1 → 0; Fig. 1.15. In this
⎨1/t if 0 < t < t1 limit, the rectangular-shaped signal occupies a re-
1
u p (t) = gion with an infinitesimally small width and a large
⎩0 if t < 0 and t > t1 . height (Fig. 1.16). The area, however, remains unity
The word ‘unit’ signifies that the signal occupies an area throughout the process. This limiting signal is known
of unity. Its Laplace transform is derived as as the unit impulse (or Dirac delta), denoted by
t1 δ(t). If the area is A = 1, it is an impulse, de-
Notation 1 −st
L[u p (t)] = Up (s) = e dt noted by Aδ(t). If an external disturbance (such as
t1 an applied force or voltage) is a pulse with very
0
large magnitude and applied for a very short period
1 − e−st1 of time, then it can be approximated as an im-
= . (1.67)
st1 pulse. Since δ(t) is the limit of u p (t) as t1 → 0, we
up (t) δ(t)
1/t1
Area = 1
Area = 1
0 0
0 t1 t 0 t
Fig. 1.14 The unit pulse u p (t) Fig. 1.16 The unit impulse δ(t)
up (t) δ(t - τ)
1/t1
Area = 1
Area = 1
0 0
0 t1 t 0 t=τ t
Fig. 1.15 The unit pulse as t1 → 0 Fig. 1.17 The unit impulse occurring at t = τ
have In general,
Part A 1.3

1 − e−st1
Notation
L[δ(t)] = Δ(s) = lim L[ f (n) (t)] = sn F(s) − sn−1 f (0) − sn−2 f˙(0)
t1 →0 st1
−st1 − · · · − f (n−1) (0) . (1.73)
L’Hôspital’s rule se
= lim =1. (1.68)
t1 →0 s
If the unit impulse occurs at some time t = τ (Fig. 1.17) Theorem 1.6: Laplace transform of integrals
it is represented by δ(t − τ), and If F(s) = L[ f (t)], then
t
L[δ(t − τ)] = e−τs . (1.69) 1
L f (τ) dτ = F(s) . (1.74)
s
This signal has the property ∞δ(t − τ) = 0 for t = τ, 0
δ(t − τ) = ∞ for t = τ, and −∞ δ(t − τ) dt = 1. It also
has the filtering property, Alternatively,
∞ t
−1 1
f (τ)δ(t − τ) dτ = f (t) . (1.70) L F(s) = f (τ) dτ . (1.75)
s
−∞ 0
1.3.3 Laplace Transform of Derivatives Solving Initial-Value Problems. The role of the
and Integrals Laplace transforms of derivatives and integrals of time-
varying functions is most significant when solving
Since engineering systems are generally modeled by an initial-value problem. Schematically, the solution
differential equations of various orders, we need to method is as in Fig. 1.18.
have knowledge of the Laplace transform of deriva-
Example 1.24: Second-order IVP
tives of different orders. In other occasions, the system
may be described by an equation that contains not only Solve ẍ + 2ẋ + x = 0, x(0) = 1, ẋ(0) = 1.
derivatives, but also integrals; for instance, a circuit in-
Solution. Laplace transformation results in
volving a resistor, an inductor, and a capacitor (RLC
circuit) [1.3]. We will also present a systematic ap- [s2 X(s) − sx(0) − ẋ(0)] + 2[sX(s) − x(0)] + X(s)
proach for solving initial-value problems.
Solve for X(s) s+3
=0 ⇒ X(s) = .
Theorem 1.5: Laplace transform of derivatives (s + 1)2
If F(s) = L[ f (t)], then Before inversion, we rewrite this last expression as
L[ f˙(t)] = sF(s) − f (0) (1.71) s+3 (s + 1) + 2
X(s) = =
and (s + 1)2 (s + 1)2
1 2
L[ f¨(t)] = s2 F(s) − s f (0) − f˙(0) . = + .
(1.72) s + 1 (s + 1)2
Time domain s domain
Algebraic equation
Initial-value problem Laplace transformation using initial conditions
in terms of
x(t) = dependent variable X(s) = transform of x(t)
Inverse Laplace transformation

x(t) Rearrange and solve for X(s)
Fig. 1.18 The solution method for initial-value problems

Finally, Case 1: Linear Factor s − pi . Each typical linear factor

Part A 1.3
s − pi of D(s) is associated with a fraction in the form

x(t) = L−1 [X(s)] = e−t + 2t e−t A
= (2t + 1) e−t . ,
s − pi
where A = const. is to be determined appropriately. We
1.3.4 Inverse Laplace Transformation
note that s = pi is called a simple pole of X(s).
Inverse Laplace transformation clearly plays a vital Example 1.25: Linear factors
role in completing the procedure for solving differen- Find L−1 [X(s)] where
tial equations. In this section we will learn a systematic s+1
technique, using partial fractions, to treat a wide range X(s) = .
(s + 3)(s + 4)
of inverse Laplace transforms. We will also introduce
the convolution method, which is quite important from Solution. The poles of X(s) are p1 = −3, p2 = −4.
a physical standpoint. Since there are two linear factors, there will be two
fractions. We write
s+1 A1 A2
Partial Fractions Method X(s) = = +
When solving an ODE in terms of x(t) through Laplace (s + 3)(s + 4) s + 3 s + 4
transformation, the very last step involves finding A1 (s + 4) + A2 (s + 3)
=
L−1 [X(s)]. And we almost always find ourselves look- (s + 3)(s + 4)
ing for the inverse Laplace transform of functions in the Collect like terms (A1 + A2 )s + 4A1 + 3A2
= ,
form of (s + 3)(s + 4)
N(s) Polynomial of degree m A1 , A2 = const.
X(s) = = , m<n. The denominators of the original and the final
D(s) Polynomial of degree n
(1.76) fractions are identical (by design), so we force their
respective numerators to be identical, that is,
The case m ≥ n is purely mathematical and does not oc- s + 1 ≡ (A1 + A2 )s + 4A1 + 3A2 .
cur in engineering analysis. The idea behind the partial
fractions method is simple: express X(s) = N(s)/D(s) But this identity holds only if the coefficients of like
as a suitable sum of fractions, and find L−1 of each powers of s on both sides are the same. So, we have
fraction accordingly. So, it all boils down to how to Coefficient of s: 1 = A1 + A2 Solve A1 = −2
break the original fraction into partial fractions, and ⇒
this depends on the nature of the poles of X(s) – the Constant term: 1 = 4A1 + 3A2 . A2 = 3
roots of D(s). In the most general sense, these roots can Insert these into the partial fractions, and perform
be real or complex. For instance, D(s) = s3 − s2 − 2s a term-by-term inverse Laplace transformation, to ob-
has roots s = 0, −1, 2, all real, so we express D(s) as tain
D(s) = s(s + 1)(s − 2). On the other hand, the roots of −2 3
X(s) = +
D(s) = s3 + 2s2 + 2s are s = 0, −1 ± i. In this case, in- s+3 s+4
stead of writing D(s) = s(s + 1 + i)(s + 1 − i) we write L−1
D(s) = s(s2 + 2s + 2). In other words, the quadratic ⇒

polynomial with complex roots remains intact. Any Linearity −1 1 −1 1
x(t) = −2L + 3L
second-degree polynomial with complex roots is called s+3 s+4
an irreducible polynomial. −3t
= −2 e + 3 e−4t
.
Structure of Partial Fractions. There are four possi-

Case 2: Repeated Linear Factor (s − pi )k . If a root
ble cases that arise in the process of finding L−1 [X(s)].
We must first identify the factors of D(s), which could of D(s), say, pi , happens to have multiplicity k, then
be linear, repeated linear, irreducible polynomial, or D(s) contains the factor (s − pi )k . This factor is then
repeated irreducible polynomial. Then, each identified associated with partial fractions
factor will have its own fraction (or fractions) associated Ak Ak−1 A2 A1
+ +· · ·+ + ,
with them. The details follow. (s − pi )k (s − pi )k−1 (s − pi )2 s − pi
where the constants Ak , · · · , A1 are determined as in Theorem 1.7: Convolution. Let G(s) = L[g(t)], H(s) =
Part A 1.3
case 1. As an example, we first write L[h(t)], and F(s) = G(s)H(s). Then,
4s + 7 L−1 [F(s)] = f (t) = (g ∗ h)(t)
X(s) =
(s + 1)2 (s + 4) t
A2 A1 A3
= + + = g(τ)h(t − τ) dτ
(s + 1)2 s + 1 s+4
Double pole; case 2 Simple pole; case 1 0
and then proceed as before to determine the constants. t

= h(τ)g(t − τ) dτ = (h ∗ g)(t) :
Case 3: Irreducible Polynomial s2 + as + b. This oc- 0
curs when X(s) has a pair of complex-conjugate poles.
The result clearly indicates that the convolution of two
Each irreducible polynomial is associated with a single
functions is symmetric.
fraction in the form
Bs + C
Example 1.26:
Convolution
,
s2 + as + b Find L−1 1
s2 (s+2)
.
where the constants B and C are found as before. Be-
fore taking the inverse Laplace transform, we must first
Solution. Write F(s) = s12 · s+2
1
= G(s)H(s) and pick
complete the square in the irreducible polynomial, that
is, s2 + as + b = (s + σ)2 + ω2 . Then, at some point, we G(s) = 1/s2 , H(s) = 1/(s + 2) so that g(t) = t, h(t) =
need to determine e−2t . Then

−1 Bs + C t
L . f (t) = (g ∗ h)(t) = τ · e−2(t−τ) dτ
(s + σ)2 + ω2
The key is to split the fraction in terms of the two ex-
0
t
pressions Integration by parts 1
= τ · e−2(t−τ)
ω s+σ 2 τ=0
and
(s + σ )2 + ω2 (s + σ)2 + ω2 t
1
so that we can ultimately use the relations [see Ta- − e−2(t−τ) dτ
ble 1.2] 2
0
ω 1 1 1
L−1 = e−σt sin ωt , = t − (1 − e−2t ) = ( e−2t + 2t − 1) .
(s + σ)2 + ω2 2 4 4

−1 s+σ
L = e−σt cos ωt . (1.77)
(s + σ)2 + ω2 1.3.5 Periodic Functions
Case 4: Repeated Irreducible Polynomial Physical systems are often subjected to external distur-
(s2 + as + b)k . The fractions are formed as bances that exhibit repeated behavior over long periods
Bk s + Ck B2 s + C2 B1 s + C1
+· · ·+ 2 + .
(s2 + as + b)k (s + as + b)2 s2 + as + b f (t)
P= 2
Convolution Method
1
In systems analysis, the problem of determining the
time history of a function often comes down to
L−1 [G(s)H(s)], where the inverse Laplace transforms
of G(s) and H(s) are known. The convolution method 0
allows us to determine L−1 [G(s)H(s)] using knowledge
of g(t) and h(t). 0 1 2 3 4 t
Notation
Notation: L−1 [G(s)H(s)] = (g ∗ h)(t) is read
“convolution of g and h” Fig. 1.19 Periodic function of Example 1.27
of time. A function f (t) is called periodic with period Solution. It is evident that the period is P = 2. With
Part A 1.4
P > 0 if it is defined for all t > 0, and f (t + P) = f (t) this, the integral in (1.78) is
for all t > 0.
2
It can then be shown [1.1] that the Laplace transform
of this function is e−st f (t) dt
0
P
1 1
F(s) = e−st f (t) dt . (1.78) f (t)=1 for 0<t<1
1 − e−Ps = e−st dt = (1 − e−s )/s .
0 f (t)=0 for 1<t<2
0
Example 1.27: Periodic signal Then, by (1.78), F(s) = (1 − e−s )/(s(1 − e−2s )). Not-
Find the Laplace transform of the periodic function in ing that1 − e−2s = 1 − ( e−s )2 = (1 − e−s )(1 + e−s ), the
Fig. 1.19. above expression reduces to F(s) = 1/(s(1 + e−s )).
1.4 Fourier Analysis

This section will focus on the concepts of Fourier P/2
series and transformation and their properties. The 2 2nπx
bn = f (x) sin dx ,
idea of Fourier series is based on representing pe- P P
riodic functions in terms of series of sine and −P/2
cosine components whose periods are integral mul- n = 1, 2, 3, · · · . (1.82)

tiples of each other. The Fourier transform is an
operation that maps a function defined in the time
domain into one in the frequency domain; its ex- In (1.79), each of the sinusoidal components with dif-
tension leads to the familiar Laplace transforma- ferent frequencies is called a harmonic, with amplitude
tion. an or bn . The harmonics basically describe the varia-
tions of f (x) about its average value of 12 a0 . In certain
1.4.1 Fourier Series cases, only a few harmonics may be needed to represent
f (x) with reasonable accuracy. On the other hand, there
Let f (x) be periodic with period P > 0, that is, f (x) are situations where many of them may be required.
is defined for all real x, and f (x + P) = f (x) for all x. At the points of discontinuity of f (x) the partial sums
Then, assuming that f (x) has a Fourier series represen- assume the average value of the left- and right-hand
tation, it is in the form limits.
∞
"
1 2nπx 2nπx
a0 + an cos + bn sin , (1.79) f (x)
2 P P
n=1
where the constants are generated by the Euler–Fourier 1

formulas, as [1.1, 7]
P/2
2
a0 = f (x) dx , (1.80)
P
−P/2 0
P/2
2 2nπx
an = f (x) cos dx ,
P P –1 0 1 x
−P/2
n = 1, 2, 3, · · · , (1.81) Fig. 1.20 A periodic function with period P = 2
Introduction to Mathematics for Mechanical Engineering 1.4 Fourier Analysis 25

Example 1.28: Fourier series Collect terms 1 2 1
Part A 1.4
Find the Fourier series representation of the periodic = − cos πx + cos 3πx + · · ·
4 π2 9
function whose description in one period is shown in
1 1 1
Fig. 1.20. + sin πx − sin 2πx + sin 3πx − · · · .
π 2 3
Solution
1 1 The third and ninth partial sums, together with the orig-
By (1.80), we have a0 = f (x) dx = x dx = 1/2. For inal function, are shown in Fig. 1.21. As mentioned
−1 0 earlier, at the points of discontinuity of f (x) the partial
n = 1, 2, 3, · · ·, (1.81) and (1.82) yield
sums assume the average value of the left- and right-
1 1 hand limits, that is, 1/2.
an = f (x) cos nπx dx = x cos nπx dx
−1 0
1.4.2 Fourier Transformation
1
= (cos nπ − 1) In Sect. 1.3 the notation F(s) was used to represent the
(nπ)2 ⎧ Laplace transform of f (t), i. e., F(s) = L[ f (t)]. Simi-
1 ⎨0 if n = even larly, fˆ(ω) is used to denote the Fourier transform of
= [(−1) n
− 1] = f (t), that is, fˆ(ω) = F[ f (t)]. Since ω is complex in gen-
(nπ)2 ⎩− 2
if n = odd ,
(nπ)2 eral, fˆ(ω) is expected to be complex-valued as well.
1 1 With this, we then write f (t) = F−1 [ fˆ(ω)] describing
bn = f (x) sin nπx dx = x sin nπx dx the inverse Fourier transform of fˆ(ω). We define the
Fourier transform pair as
−1 0
1 1 ∞
1
= − [x cos nπ]10 + [sin nπx]10 ˆf (ω) = √ f (τ) e−iωτ dτ , (1.83)
nπ (nπ)2 2π
1 −∞
= (−1)n+1 . ∞
nπ 1
Equation (1.79) gives the Fourier series of f (x), as f (t) = √ fˆ(ω) eiωt dω . (1.84)
2π
1 2 1 1 −∞
− 2 cos πx + sin πx − sin 2πx Fourier transformation can be thought of as a map-
4 π π 2π
2 1 ping that assigns to a given function of time t an
− 2 cos 3πx + sin 3πx + · · · integral function of frequency ω. In general, any trans-
9π 3π
f (x)
1.2
1.0
9th partial sum
0.8
0.6
0.4
3rd partial sum
0.2
– 0.2
–1.5 – 1.0 – 0.5 0 0.5 1.0 1.5
x
Fig. 1.21 Example 1.28
Table 1.3 Fourier transform pairs formation with this type of property is known as an
Part A 1.5
integral transformation. The obvious similarities be-

No. f (t) f̂ (ω)
⎧ tween the Laplace and Fourier transforms are credited
⎨1 −b < t < b to the Laplace transformation being an integral trans-
2 sin bω
1 π ω
⎩0 otherwise formation itself. Fourier transforms of several functions
⎧ are listed in Table 1.3.
⎨1 b < t < b −ib1 ω − e−ib2 ω
1 2
2 √1 e
⎩0 otherwise 2π iω Example 1.29: Fourier transform
⎧ Find the Fourier transform of
⎨ e−at t > 0 ⎧
3 ,a > 0 √1 1 ⎨0 if t < 0
⎩0 2π a+iω
otherwise f (t) = (a > 0) .
⎧ ⎩ e−at if t > 0
⎨ eat t < 0
4 ,a > 0 √1 1
⎩0 otherwise
2π a−iω
⎧
⎨ eat b < t < b Solution. By (1.83),
1 2 (a−iω)b2 − e(a−iω)b1
5 √1 e
⎩0 otherwise
2π a−iω
∞
1
6 e−a|t| , a > 0 2 a fˆ(ω) = √ f (τ) e−iωτ dτ
π ω2 +a2 2π
⎧ −∞
⎨− e−at t < 0
7 ,a < 0 2 −iω ∞
⎩ eat π ω2 +a2 1
t>0 =√ e−aτ e−iωτ dτ
⎧ 2π
⎨ eiat −b < t < b 0
8
⎩0
2 sin(ω−a)b
π ω−a 1 −1 −(a+iω)τ ∞
otherwise =√ e
⎧ 2π a + iω 0
⎨ eiat b < t < b i(a−ω)b1 − ei(a−ω)b2 1 1
9
1 2
√i e =√ .
⎩0 otherwise
2π a−ω
2π a + iω
π e−a|ω|
1
,a > 0
10
a2 +t 2
⎧
2 a Using fˆ(ω) above in (1.84), we find
⎨t 0 < t < b −ibω (1+ibω)
11 √1 −1+ e ∞
⎩0 otherwise 2π ω2
1 1 1
⎧ f (t) = √ √ eiωt dω
⎪
⎪ 0<t<b 2π 2π a + iω
⎨t ibω − e−2ibω
−∞
12 2t − b b < t < 2b √1 −1+2 e ∞
⎪
⎪ 2π ω2
1 1
⎩
0 otherwise = eiωt dω .
2π a + iω
13
2
e−at , a > 0
2
√1 e−ω /(4a) −∞
2a
2 /(4a) √ −aω2
14 e−t ,a > 0 2a e This is known as the complex Fourier integral represen-
tation of the function under consideration.
1.5 Linear Algebra

In this section we present the fundamentals of linear equations in the system involve many of the unknown
algebra, specifically, vectors and matrices, and their re- variables. In these cases, techniques such as direct sub-
lation to linear systems of algebraic and differential stitution and elimination are no longer suitable due to
equations. The methods of linear algebra are mainly their lack of computational efficiency. We focus on al-
useful in the treatment of systems of equations that gebraic systems first, then extend the ideas to systems
are heavily coupled, that is, when a large number of of differential equations.
Introduction to Mathematics for Mechanical Engineering 1.5 Linear Algebra 27
1.5.1 Vectors and Matrices of A. These diagonal elements form the main diagonal
Part A 1.5
of A. The diagonal directly below the main is known
An n-dimensional vector v is an ordered set of n as the subdiagonal and the one above the main is called
scalars, and is written as v = (v1 , v2 , · · · , vn ). Each vi the superdiagonal. Two matrices A = [aij ] and B = [bij ]
(i = 1, 2, · · · , n) is called a component of the vector v. are said to be equal if they have the same size, and the
In a general sense, an n-dimensional vector helps us lo- same entries in the respective locations. If some rows
cate a point in an n-dimensional space, regardless of or columns (or possibly both) of A are deleted, the out-
the nature of its components or perhaps what physical come is a submatrix of A. If no rows or columns of A are
quantity each may represent. omitted, we have A as a submatrix of itself. Submatri-
ces play important roles in such areas of matrix analysis
Matrices as determinants and rank.
A collection of numbers (real or complex) or possibly
functions, arranged in a rectangular array and enclosed Matrix Operations
by brackets, is referred to as a matrix. Each of the ele- Matrices of the same size can be added. The result,
ments in a matrix is called an entry (or element) of the or the sum, is a matrix of the same size. If A = [aij ]
matrix. The horizontal and vertical lines are referred to and B = [bij ] are m × n, their sum C = [cij ] is also
as rows and columns of the matrix, respectively. A ma- m × n. Matrix addition is performed entry-wise, that is,
trix is called a row vector if it consists of one row only, the entry of C in the (i, j) slot is the sum of the entries
and a column vector if it has one column only. The num- of A and B in that same slot. The m × n zero matrix, de-
ber of rows and columns of a matrix determine the size noted by 0m×n , is an m × n matrix all of whose entries
of that matrix. If a matrix A has m rows and n columns, are zero. If A = [aij ] is m × n and k is a scalar, then kA
then it is said to be of size m × n. If the number of rows is an m × n matrix whose entries are those of A mul-
and columns are the same, we speak of a square matrix, tiplied by k in every slot, that is, kA = [kaij ]m×n . Let
otherwise, a rectangular matrix. We denote matrices by A = [aij ]m×n and B = [bij ]n× p . It is important to note
bold-faced capital letters, such as A. The abbreviated that the number of columns of A is n, which is equal to
form of an m × n matrix is the number of rows of B. Then, their product C = AB
is m × p whose entries are obtained as
A = [aij ]m×n , i = 1, 2, · · · , m , j = 1, 2, · · · , n ,
where aij is known as the (i, j) entry of A, located "
n
cij = aik bk j , i = 1, 2, · · · , m , j = 1, 2, . . . , p .
at the intersection of the ith row and the jth column
k=1
of A so that a12 , for instance, occupies the entry at
(1.85)
which the first row and the second column meet. In the
event that A is a square matrix (m = n), the elements This is shown schematically in Fig. 1.22. If the number
a11 , a22 , · · · ,ann are referred to as the diagonal entries of columns of A does not match the number of rows of
a11 a12 a1n c11 c12 c1j c1p

jth column
a21 a22 a2n c21 c22 c2 j c2 p
b11 b12 b1j b1p
b21 b22 b2 j b2 p
=
ith row ai1 ai 2 ain ci1 ci2 cij cip
(i, j ) entry
bn1 bn2 bnj bnp
am1 am2 amn cm1 cm2 cmj cmp
Fig. 1.22 Construction of the matrix product AB = C

B, the product is undefined. If the product is defined, Example 1.31: Special matrices
Part A 1.5
then to get the (i, j) entry of C, we proceed as follows: Matrices U, L, and D are upper triangular, lower trian-
the ith row of A is clearly a 1 × n vector. The jth column gular, and diagonal, respectively:
of B is an n × 1 vector, hence these two vectors have the ⎛ ⎞ ⎛ ⎞
same number of components, n. In these two vectors, −2 1 2 1 0 0
⎜ ⎟ ⎜ ⎟
multiply the first components, the second components, U=⎝ 0 5 0⎠ , L=⎝2 0 0⎠ ,
etc., up to the nth components. Then add the individual 0 0 3 4 7 −1
products together. The result is cij . ⎛ ⎞
3 0 0
⎜ ⎟
D = ⎝ 0 −4 0 ⎠ .
Example 1.30: Matrix Multiplication
Find 0 0 1
⎛
⎞
# $ −2 −1 4 Note that in U and L zeros are allowed along the main
1 −2 3 ⎜ ⎟ diagonal. In fact, the main diagonal may consist of all
AB = ⎝ 1 2 0⎠ .
0 1 4 2×3 zeros. On the other hand, D may have one or more zero
3 5 1 3×3 diagonal elements, as long as they are not all zeros. In
the event that all entries of an n × n matrix are zeros, it
Solution. We first note that the operation is valid be- is called the n × n zero matrix 0n×n .
cause A has three columns and B has three rows. And,
AB will be 2 × 3. Following the strategy outlined above,
Determinant
we find the product as
The determinant of a square matrix A = [aij ]n×n is a real

scalar denoted by |A| or det(A). For the most trivial case
AB = 1·(−2)+(−2)·1+3·3 1·(−1)+(−2)·2+3·5
0·(−2)+1·1+4·3

0·(−1)+1·2+4·5 of n = 1, A = [a11 ], and we define the determinant sim-
1·4+(−2)·0+3·1 ply as |A| = a11 . For n ≥ 2, the determinant is defined
0·4+1·0+4·1
# $ as
5 10 7 using the i-th row
= .
13 22 4 2×3 "
n
|A| = aik (−1)i+k Mik , i = 1, 2, · · · , n (1.89)
k=1
Matrix Transpose
or
Given an m × n matrix A, its transpose, denoted by AT ,
using the j-th column
is an n × m matrix with the property that its first row
is the first column of A, its second row is the second "
n
column of A, and so on. Given that all matrix operations |A| = ak j (−1)k+ j Mk j , j = 1, 2, · · · , n (1.90)
are valid, k=1
Here Mik is the minor of the entry aik , defined as the
(A + B)T = AT + BT (1.86)
determinant of the (n − 1) × (n − 1) submatrix of A ob-
(kA) = kA ,
T T
scalar k (1.87) tained by deleting the ith row and the kth column of A.
(AB)T = BT AT . (1.88) The quantity (−1)i+k Mik is known as the cofactor of
aik and is denoted by Cik . Also note that (−1)i+k is re-
sponsible for whether a term is multiplied by +1 or −1.
Special Matrices Equations (1.89) and (1.90) suggest that the determinant
A square matrix A is symmetric if AT = A and skew- of a square matrix can be calculated using any row or
symmetric if AT = −A. A square matrix An×n = [aij ] any column of the matrix. However, for all practical pur-
is called upper-triangular if aij = 0 for all i > j, that poses, it is wise to use the row (or column) containing
is, every entry below the main diagonal is zero, lower- the most number of zeros, or if none, the one with the
triangular if aij = 0 for all i < j, that is, all elements smallest entries. A square matrix with a nonzero deter-
above the main diagonal are zeros, and diagonal if minant is known as a nonsingular matrix. Otherwise, it
aij = 0 for all i = j. The n × n identity matrix is a di- is called singular. The rank of any matrix A, denoted by
agonal matrix whose diagonal entries are all equal to 1, rank(A), is the size of the largest nonsingular submatrix
and is denoted by I. of A. If |An×n | = 0, we conclude that rank (A) < n.
Example 1.32: 3 × 3 determinant Determinant of Block Matrices. We define a block-
Part A 1.5
Find the determinant of diagonal matrix as a square matrix partitioned such
⎛ ⎞
1 2 −3 that its diagonal elements are square matrices, while
⎜ ⎟ all other elements are zeros; see Fig. 1.23a. Similarly,
A = ⎝ 4 −1 1 ⎠ .
a block-triangular matrix is a square matrix partitioned
2 0 1 so that its diagonal elements are square blocks, while all
entries either above or below this main block diagonal
are zeros; see Fig. 1.23b,c.
Solution. We will use the third row because it happens
Many properties of these special block matrices are
to contain a zero. Following (1.89),
basically extensions of those of diagonal and triangu-
|A| = 2 · (−1)3+1 M31 + 0 + 1 · (−1)3+3 M33 lar matrices. In particular, the determinant of each of
these matrices is equal to the product of the individ-

2 −3 1 2 ual determinants of the blocks along the main diagonal.
= 2 +
−1 1 4 −1 Consequently, a block diagonal (or triangular) matrix is
singular if and only if one of the blocks along the main
= 2(2 − 3) + (−1 − 8) = −11 .
diagonal is singular.
Properties of Determinant. The determinant of a ma- Inverse of a Matrix. Given a square matrix An×n , its
trix possesses a number of important properties, some inverse is denoted by A−1 with the property that
of which are listed below [1.1]:
AA−1 = I = A−1 A , (1.91)
• A square matrix A and its transpose
have the same where I denotes the n × n identity matrix. If A−1 exists,
determinant, that is, |A| = AT . then it is unique. A square matrix has an inverse if and
• The determinant of diagonal, upper-triangular and only if it is nonsingular. Equivalently, An×n has an in-
lower-triangular matrices is the product of the diag- verse if and only if rank (A) = n. A square matrix with
onal entries. an inverse is called invertible. An immediate applica-
• If an entire row (or column) of a square matrix A is tion of the inverse is in the solution process of a linear
zero, then |A| = 0. system Ax = b. Multiplying this equation from the left,
• If A is n × n and k is scalar, then |kA| = kn |A|. known as premultiplication, by A−1 , yields
• If any two rows (or columns) of A are interchanged,
the determinant of the resulting matrix is − |A|. A−1 (Ax) = A−1 b
• The determinant of the product of two matrices ⇒ (A−1 A)x = A−1 b
obeys |AB| = |A| |B|. ⇒ Ix = A−1 b
• Any square matrix with any number of linearly de-
pendent rows (or columns) is singular. ⇒ x = A−1 b .
a) b) c)
0 0
*
0 0 *
Fig. 1.23 (a) Block-diagonal matrix. (b) Block-upper-triangular matrix. (c) Block-lower-triangular matrix
Inverse via the Adjoint Matrix. The inverse of an invert- the original matrix. The inverse of an upper-triangular
Part A 1.5
ible matrix A = [aij ]n×n is determined using the adjoint matrix is upper-triangular. The diagonal elements of
of A, denoted by adj(A) and defined as [1.1] the inverse are the reciprocals of the diagonal entries
of the original matrix, while the off-diagonal entries
adj(A) do not obey any pattern. A similar result holds for
⎛ ⎞
(−1)1+1 M11 (−1)2+1 M21 · · · (−1)n+1 Mn1 lower-triangular matrices. Furthermore, it turns out that
⎜ ⎟a block-diagonal matrix and its inverse have exactly the
⎜ (−1)1+2 M12 (−1)2+2 M22 · · · (−1)n+2 Mn2 ⎟
=⎜⎜ .. .. .. ⎟
⎟
same structure.
⎝ . . . ⎠
(−1)1+n M1n (−1)2+n M2n · · · (−1)n+n Mnn Properties of Inverse. Some important properties of the
⎛ ⎞ inverse [1.1, 8] are given below. The assumption is that
C11 C21 · · · Cn1 all listed inverses exist.
⎜ ⎟
⎜ C12 C22 · · · Cn2 ⎟
(1.92) • (A ) = A.
−1 −1
=⎜⎜ .. .. .. ⎟
⎟.
⎝ . . . ⎠ • (AB)−1 = B−1 A−1 .
C1n C2n · · · Cnn • (AT )−1 = (A−1 )T .
• The inverse of a symmetric matrix is symmetric.
Note that each minor Mij (or cofactor Cij ) occupies the • (A p )−1 = (A−1 ) p , where p is a positive integer.
( j, i) position in the adjoint matrix, the opposite of what • det(A−1 ) = 1/ det(A).
one would normally expect. Then, the inverse of A is
simply defined by
1.5.2 Eigenvalues and Eigenvectors
−1 1
A = adj(A) . (1.93)
|A| The fundamentals of linear algebra are now extended
to treat systems of differential equations, which are
of particular importance to us since they represent the
Example 1.33: Formula for the inverse of a 2 × 2 matrix
mathematical models of dynamic systems. In the anal-
Find a formula for the inverse of ysis of such systems, one frequently encounters the
# $
a11 a12 eigenvalue problem, solutions of which are eigenvalues
A= . and eigenvectors. This knowledge enables the analyst to
a21 a22
determine the natural frequencies and responses of sys-
tems. Let A be an n × n matrix, v a nonzero n × 1 vector,
Solution. Following the procedure outlined above, we and λ a number (complex in general). Consider
find Av = λv (1.95)
M11 = a22 , C11 = a22 , A number λ for which (1.95) has a nontrivial solution
M12 = a21 , C12 = −a21 , (v = 0n×1 ) is called an eigenvalue or characteristic value
M21 = a12 , C21 = −a12 , of matrix A. The corresponding solution v = 0 of (1.95)
is the eigenvector or characteristic vector of A corre-
M22 = a11 , C22 = a11 . sponding to λ. Eigenvalues, together with eigenvectors
Then, form the eigensystem of A. The problem of determin-
# $ ing eigenvalues and the corresponding eigenvectors of
−1 1 a22 −a12 A, described by (1.95), is called an eigenvalue problem.
A = , (1.94) The trace of a square matrix A = [aij ]n×n , denoted by
|A| −a21 a11
tr(A), is defined as the sum of the eigenvalues of A. It
which is a useful formula for 2 × 2 matrices, allowing us turns out that tr(A) is also the sum of the diagonal el-
to omit the intermediate steps. ements of A. A matrix and its transpose have the same
eigenvalues.
Inverses of Special Matrices. If the main diagonal en-
tries are all nonzero, the inverse of a diagonal matrix Solving the Eigenvalue Problem
is again diagonal. The diagonal elements of the inverse Let us consider (1.95), Av = λv. Because equations in
are simply the reciprocals of the diagonal elements of this form involve scalars, vectors, and matrices, it is im-
perative that extra caution is taken while working with we apply suitable elementary row operations [1.1] to the
Part A 1.5
them. First, rewrite and manipulate (1.95) as augmented
# matrix
$ to reduce it to
Av − λv = 0n×1 ⇒ (A − λI)v = 0 , 0 1 0
(1.96) .
0 0 0
where we note that every term here is an n × 1 vector. The second row suggests that there is a free variable,
The identity matrix I = In has been inserted so that the implying that the two equations contained in (1.99) are
two terms in parentheses are compatible; otherwise we linearly dependent. From the first row, we have v21 =
would have A − λ, which is meaningless. This equa- 0 so that v21 cannot be the free variable, so v11 must
tion has a nontrivial solution (v = 0) if and only if the be. In this example, since we already have v21 = 0, then
coefficient matrix, A − λI, is singular. That means v11 = 0 because otherwise v1 = 0, which is not valid.
|A − λI| = 0 . (1.97) # $ let v11 = 1, so
For simplicity,
1
This is called the characteristic equation of A. The de- v1 = .
0
terminant |A − λI| is an nth-degree polynomial in λ
Similarly, the eigenvector corresponding to λ2 = 2 can
and is known as the characteristic polynomial of A
be shown to be v2 = [−1 1]T . The set (v1 , v2 ) is the
whose roots are precisely the eigenvalues of A. Once
basis of all eigenvectors of matrix A.
the eigenvalues have been identified, each eigenvector
corresponding to each of the eigenvalues is determined Special Matrices
by solving (1.96). The eigenvalues of triangular and diagonal matrices
Example 1.34: Eigenvalues and eigenvectors are the diagonal entries. The eigenvalues of block-
Find the eigenvalues and eigenvectors of triangular and diagonal matrices are the eigenvalues of
# $ the block matrices along the main diagonal. All eigen-
−1 −3 values of a symmetric matrix are real, while those of
A= . a skew-symmetric matrix are either zero or pure imagi-
0 2
nary.
Solution. To find the eigenvalues of A, we solve the Generalized Eigenvectors
characteristic equation, If λk is an eigenvalue of A occurring m k times,
then m k is the algebraic multiplicity of λk , denoted
|A − λI| = 0
by AM(λk ). The maximum number of linearly in-
dependent eigenvectors associated with λk is called
−1 − λ −3
⇒ =0 the geometric multiplicity of λk , GM(λk ). In gen-
0 2−λ
eral, GM(λk ) ≤ AM(λk ). In Example 1.34 the AM and
⇒ (λ + 1)(λ − 2) = 0 GM of each of the two eigenvalues was 1. When
⇒ λ1,2 = −1, 2 . GM(λk ) <AM(λk ), there are fewer eigenvectors than
one would expect. For instance, if AM(λ) = 2 and
Without losing any information, let us assign λ1 = −1. GM(λ) = 1, then only one independent eigenvector can
To find the eigenvector, solve (1.96) with λ = λ1 = −1, be found for λ, while one is missing; the missing one is
λ =−1 called a generalized eigenvector [1.1].
(A − λ1 I)v1 = 0 ⇒ (A + I)v1 = 0 ,
1
(1.98)
Similarity Transformation – Diagonalization
where v1 is the
# 2 ×$1 eigenvector corresponding to λ1 . Two matrices An×n and Bn×n are said to be sim-
v11 ilar if there exists a nonsingular matrix Sn×n such
Letting v1 = and using A in (1.98), we find
v21 that B = S−1 AS. We say that B is obtained from
# $# $ # $ A through a similarity transformation. Similar matri-
0 −3 v11 0 ces have the same eigenvalues. Suppose An×n has n
= . (1.99)
linearly independent eigenvectors v1 , v2 , · · · ,vn asso-
0 3 v21 0
ciated with eigenvalues λ1 , λ2 , · · · , λn . Form the n × n
As expected, this system has nontrivial solutions be- matrix P = [ v1 v2 . . vn ], known as the modal ma-
cause the coefficient matrix is singular. To solve (1.99), trix, whose columns are the eigenvectors of A. Then, P
is nonsingular, and 2. What are chosen as the state variables?

Part A 1.5
⎛ ⎞ Those variables for which initial conditions are re-

λ1 0 quired in (1) are chosen as the state variables.
⎜ ⎟
⎜ · ⎟
P−1 AP = = ⎜ ⎟, (1.100)
⎝ · ⎠ Example 1.36: State variables
The equation of motion for the mass–spring–damper
0 λn
system in Fig. 1.24 is given as
where the order of the appearance of the λi along the
m ẍ + cẋ + kx = f (t) .
main diagonal agrees with the order of the correspond-
ing vi in matrix P. Any matrix A that satisfies (1.100) For a complete description, we need the initial con-
is called diagonalizable. It is clear that A and are ditions x(0) = x0 (initial displacement) and ẋ(0) = v0
similar, hence share the same eigenvalues. (initial velocity).
Example 1.35: Modal matrix
Since two initial conditions are required, there are
In reference to Example 1.34, the modal matrix is two state variables. And, since these conditions are re-
# $ quired for x and ẋ, the state variables are
1 −1 x1 = x , x2 = ẋ .
P= , which satisfies
0 1
# $ # $
−1 0 λ1 0 State-Variable Equations. The idea now is to use the
−1
P AP = = = . selected set of state variables to transform the original
0 2 0 λ2
system into a larger system of first-order ODEs. Each of
the resulting differential equations consists of the time
derivative of one of the state variables on one side, and
1.5.3 Numerical Solution a linear combination of the state variables and possibly
of Higher-Order Systems of ODEs system inputs on the other side. Each of these first-order
ODEs just described is called a state-variable equation.
In this section we present a two-step process to treat one
or more ODEs of order higher than one. The first stage Example 1.37: State-variable equations
is concerned with the conversion of the original ODEs Referring to Example 1.36, since there are two state
into a system of first-order ODEs. The second stage variables, two first-order ODEs must be obtained in
deals with the application of the fourth-order Runge– terms of x1 and x2 . The first equation is simply ẋ1 = x2 .
Kutta method to solve this system numerically; also To obtain the second equation, we use the equation of
see Sect. 1.2.2. It should be mentioned that we are not motion as
interested in mathematically fabricated, higher-order 1
ẋ2 = [−kx1 − cx2 + f (t)] .
systems, but rather those of physical significance. The m
first step is handled via the state variables, explained The state-variable equations are then
below. ⎧
⎨ẋ = x
1 2
State Variables ⎩ẋ = 1 [−kx − cx + f (t)] .
2 m 1 2
The smallest possible set of independent variables that
completely describes the state of a system is referred These can be expressed in matrix form, as
to as the set of state variables. Since independence is ẋ = f (t, x) , x(0) = x0 , (1.101)
essential, state variables cannot be expressible as alge-
braic functions of one another. Moreover, the set of state where

variables for a certain system is not unique. Given a set x1
of differential equations describing a certain physical x= ,
x2
system, two key questions need be answered [1.1, 3].
x2
f (t, x) = ,
1. How many state variables are there? −(k/m)x1 − (c/m)x2 + (1/m) f (t)
The number of state variables is equal to the num-
x0
ber of initial conditions required to completely solve the x0 = .
system’s governing equations. v0
Introduction to Mathematics for Mechanical Engineering References 33
Fourth-Order Runge–Kutta Method for Systems
Part A 1
Displacement
Numerical solution of the state-variable equations – x
such as that in (1.101) of Example 1.37 – is then ob- Spring k
tained via the extension of RK4 discussed in Sect. 1.2.2.
Consider a system in the form f (t)
m
ẋ(t) = f (t, x(t)) , x(a) = x0 , a≤t≤b, Applied force
(1.102)
Damper c Mass
where
⎛ ⎞ Fig. 1.24 A mechanical system
x1 (t)
⎜ ⎟
⎜ x2 (t) ⎟ N subintervals. The fourth-order Runge–Kutta method
x(t) = ⎜
⎜ .. ⎟ ,
⎟
⎝ . ⎠ (RK4) for a system of first-order ODEs is as fol-
lows [1.5]. Knowing the initial vector x0 , the solution
xn (t) vector xi at each of the subsequent mesh points ti is
⎛ ⎞
f 1 (t, x1 , x2 , · · · , xn ) obtained via
⎜ ⎟
⎜ f 2 (t, x1 , x2 , · · · , xn ) ⎟ 1
⎜
f (t, x(t)) = ⎜ .. ⎟, xi+1 = xi + [q1 + 2q2 + 2q3 + q4 ],
⎟ 6
⎝ . ⎠
i = 0, 1, 2, · · · , N − 1 ,
f n (t, x1 , x2 , · · · , xn )
⎛ ⎞ where
α1
⎜ ⎟
⎜ α2 ⎟ q1 = h f (ti , xi ) ,
⎜
x0 = x(a) = ⎜ . ⎟ ⎟.

⎝ .. ⎠ 1
q2 = h f ti + h, xi + q1 ,
1
2 2
αn
1 1
Define an integer N > 0 and let h = (b − a)/N be q3 = h f ti + h, xi + q2 ,
the step size. The mesh points ti = a + ih, i = 2 2
0, 1, · · · , N − 1, then partition the interval [a, b] into q4 = h f (ti + h, xi + q3 ) .
References
1.1 R.S. Esfandiari: Applied Mathematics for En- 1.5 J.H. Mathews: Numerical Methods for Computer Sci-
gineers, 4th edn. (Atlantis, Irvine, California ence, Engineering, and Mathematics (Prentice-Hall,
2008) New York 1987)
1.2 J.W. Brown, R.V. Churchill: Complex Variables and 1.6 R.L. Burden, J.D. Faires: Numerical Analysis, 3rd edn.
Applications, 7th edn. (McGraw-Hill, New York (Prindle, Boston 1985)
2003) 1.7 J.W. Brown, R.V. Churchill: Fourier Series and
1.3 H.V. Vu, R.S. Esfandiari: Dynamic Systems: Modeling Boundary Value Problems. 6th edn. (McGraw-Hill,
and Analysis (McGraw-Hill, New York 1997) New York 2001)
1.4 C.H. Edwards, D.E. Penney: Elementary Differential 1.8 G.H. Golub, C.F. Van Loan: Matrix Computations, 3rd
Equations with Boundary Value Problems, 4th edn. edn. (The Johns Hopkins University Press, London
(Prentice-Hall, New York 2000) 1996)

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1.1 Complex Analysis .................................. 4

1.1 Complex Analysis

i 3+i Geometrically, |z| is the distance from z to the origin of

Division of Complex Numbers

Complex conjugation is extremely useful in com-

bra, in particular, multiplication and division. Consider

z 1 z 2 = r1r2 ei(θ1 +θ2 ) or r1 r2 (θ1 + θ2 ) . (1.15) r

|z 1 z 2 | = r1r2 = |z 1 ||z 2 | and –θ

Similarly, for division of complex numbers, we have z = r e -iθ

neighborhood of z 0 . A function is analytic in a domain

Test of Analyticity: Cauchy–Riemann Equations.

w = f (z) = u(x, y) + iv(x, y) . (1.22) Partial differentiation yields

1.2 Differential Equations

where y = y(x) and y(n) = dn y/ dx n . If all coefficients Example 1.8: IVP

Equations wise it is called nonhomogeneous.

The functions y1 = e2x and y2 = e−3x are linearly in-

where Nonhomogeneous Linear Second-Order ODEs

Solution. We first find the characteristic equation and yh = c1 y1 + c2 y2 , (c1 , c2 constants)

x +1 and e−x . Method of Undetermined Coefficients. The idea intro-

Example 1.17: Special case II Step 3: General Solution. Combination of yh and yp

1.3 Laplace Transformation

L[a1 f 1 (t) + a2 f 2 (t)]

16 sin ωt ω Example 1.21: Shift on the s-axis

22 t cos ωt s 2 −ω2 of Laplace Transforms

ω21 = ω22 F (s + a) F (s)

t f (t) = −L−1 [F  (s)] . (1.56)

Theorem 1.4: Shift on the t-axis. Given that F(s) =

Finding L[u(t − a)] via Theorem 1.4. We now have the

nal is called a ramp function, denoted by Au r (t). In that Au p (t), and

Time domain s domain

Inverse Laplace transformation

Fig. 1.18 The solution method for initial-value problems

Finally, Case 1: Linear Factor s − pi . Each typical linear factor

s − pi of D(s) is associated with a fraction in the form

Structure of Partial Fractions. There are four possi-

and then proceed as before to determine the constants. t

1.4 Fourier Analysis

cosine components whose periods are integral mul- n = 1, 2, 3, · · · . (1.82)

where the constants are generated by the Euler–Fourier 1

integral transformation. The obvious similarities be-

1.5 Linear Algebra

a11 a12 a1n c11 c12 c1j c1p

am1 am2 amn cm1 cm2 cmj cmp

Fig. 1.22 Construction of the matrix product AB = C

Example 1.32: 3 × 3 determinant Determinant of Block Matrices. We define a block-

is nonsingular, and 2. What are chosen as the state variables?

⎛ ⎞ Those variables for which initial conditions are re-

Fourth-Order Runge–Kutta Method for Systems

S-ar putea să vă placă și

t f (t) = −L−1 [F (s)] . (1.56)